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warning/0506/astro-ph0506138.html	ar5iv	text	# Spectral energy distribution of the metagalactic ionizing radiation field from QSO absorption spectraBased on observations obtained at the VLT Kueyen telescope (ESO, Paranal, Chile), the ESO programme 65.O-0474(A) ## 1 Introduction The study of the spectral energy distribution (SED) in the radiation background resulting from flux emitted by all celestial sources is an important part of modern observational cosmology. The current paradigm assumes that the metagalactic ionizing background is formed by radiation of QSOs and galaxies (stars) reprocessed by the intergalactic medium (IGM). The SED evolves with cosmic time due to different contribution of galactic and QSO radiation, and to varying IGM opacity caused by the hydrogen and helium reionization. Since the pioneering works of Chaffee et al. (1986) and Bergeron & Stasinska (1986), the approach to recover the shape of the ionizing radiation field from the measurements of the intervening metal absorbers in QSO spectra has been employed in numerous studies. In particular, the energy range 1 Ryd $`<E<`$ 10 Ryd can be probed through the relative intensities of metal lines such as Si ii-Si iv, C ii-C iv, N iii, N v, O vi. A common procedure consists of selecting a standard SED and checking whether it is consistent with measured column densities. This procedure can be appropriate for obtaining some general information about the SED, but its effectiveness and accuracy is rather low because of lack of search strategy. In the present paper we describe a computational technique exploiting the response surface methodology from the theory of experimental design which enables a directed search for the shape of the ionizing background. The proposed method reliably recovers the main features in the spectral shape of the metagalactic flux by the analysis of optically thin QSO absorption systems. Four absorption systems were selected to illustrate how this approach can be implemented in practice. Three of these systems have redshifts $`z3`$. This redshift is of particular interest for studying the traces of still not completely ionized He ii (Reimers et al. 1997). In general, the presence of He ii in the intergalactic medium affects the ionizing spectrum in the range $`E>1`$ Ryd due to He ii Ly$`\alpha `$, two-photon reemission and to He ii continuum absorption (Haardt & Madau 1996; Fardal et al. 1998). The He ii Ly$`\alpha `$ absorption troughs detected in spectra of five quasars (Vogel & Reimers 1995; Zheng et al. 1998; Anderson et al. 1999; Heap et al. 2000; Smette et al. 2002; Jakobsen et al. 2003 ; Zheng et al. 2004a) may indicate He ii Gunn-Peterson effect (He ii Ly$`\alpha `$ absorption in a diffuse IGM) at $`z3`$. Furthermore, recent FUSE observations of the He ii Ly$`\alpha `$ forest towards HE 2347–4342 combined with observations of the H i Ly$`\alpha `$ forest at the VLT revealed also large scale variations in $`\eta =He\text{ii}/H\text{i}`$ which are still not well understood (Shull et al. 2004; Zheng et al. 2004b). One possible explanation is the presence of spatial fluctuations in the metagalactic ionizing field at $`E>\mathrm{\hspace{0.33em}4}`$ Ryd. To confirm this suggestion, additional observations at different redshifts and in different sightlines are needed. However, direct measurements of He ii Ly$`\alpha `$ absorption at 303.78 Å are very problematic due to almost complete light blotting of distant QSOs by the intervening Lyman limit or damped Ly$`\alpha `$ systems (LLSs and DLAs, respectively). In this context, the proposed approach to estimate the shape of the ionizing background from metal absorption systems is of great importance. The structure of the paper is as follows. The procedure to recover the shape of the underlying UV continuum and an example illustrating how to use it are described in Sect. 2. This section also contains a brief description of a computational method used to invert the observed line profiles – the Monte Carlo Inversion (MCI). The detailed analysis of absorption systems used for the SED estimations is given in Sect. 3. The results are discussed and summarized in Sect. 4. An example of an experimental design is given in the Appendix. ## 2 Computational methods ### 2.1 Shape of the ionizing radiation Here follows a description of how the SED of the ionizing radiation can be estimated from metal line profiles observed in optically thin absorption systems. The method is based on the response surface methodology used in experimental design (see, e.g., Box, Hunter & Hunter 1978, Chapter 15). A basic UV spectrum is taken as an initial guess. The estimated column densities of different ions and a photoionization code (e.g., CLOUDY, Ferland 1997) are used to derive the ionization parameter, metallicity and element abundance ratios. If the observed column densities are well reproduced with this trial spectrum, it implies that the initial guess was appropriate. If, however, the observed column densities cannot be reproduced or/and some other inconsistencies arise, then the assumed shape of the UV background is to be adjusted. To allow for quantitative estimations, the shape of the UV continuum must be parameterized. In general, the shape of the ionizing spectrum can be specified by a piecewise continuous function, e.g. by a set of power laws and/or exponents. In particular, in the range 1 Ryd $`<E<10`$ Ryd relevant for the ions frequently observed in QSO spectra, the following variables (called ‘factors’ in the experimental design) can be used to describe the main features of the SED (Fig. 1): A first region between 1 Ryd and the He i break is characterized by a slope $`f_1`$ (power law exponent) and by the position $`f_2`$ of the He i break. A second region can be defined between the He i and He ii breaks (point $`A`$) with $`f_3`$ and $`f_4`$ the slope and the position of the He ii break, respectively. Since the energy depression can be relevant, we introduce here also $`f_5`$ which defines the depth of the break \[$`f_5=\mathrm{log}(J_\mathrm{B}/J_\mathrm{A})`$\], and $`f_6`$ – the slope between points $`A`$ and $`B`$. $`f_7`$ characterizes the energy spectrum in the region $`BC`$. To account for a possible bump around 3 Ryd due to recombinations within the clumpy intergalactic gas (He ii Ly$`\alpha `$ emission, He ii two-photon continuum emission and He ii Balmer continuum emission), we introduce factors $`f_8`$ and $`f_9`$ which describe the amplitude of the bump and its width, respectively. The energy of the far UV cut off (point $`C`$) when taken above 100 Ryd does not affect the fractional ionizations of ions we are interested in. In all computations described in subsequent sections this energy and the slope after it are kept fixed at 128 Ryd and –1.5, respectively. Note that the positions of the He i and He ii breaks can be shifted due to large scale motions and/or superposition of different spectra (e.g. stellar+metagalactic), and hence, in general the values of $`f_2`$ and $`f_4`$ may not be equal to 1.8 Ryd and 4 Ryd, respectively. Additional factors can be defined to describe some particular features of the SED as, for example, those caused by He ii Ly$`\alpha `$ absorption in a diffuse IGM (see Sect. 3 below). In this notation, an ionizing spectrum is determined as a point $`\{f_1,\mathrm{},f_k\}`$ in the $`k`$-dimensional factor space. The number and selection of factors depend on a particular spectral shape. For instance, the metagalactic spectrum at $`z=3`$ calculated by Haardt & Madau (1996, hereafter HM) has coordinates $`\{f_1,f_4,f_5,f_6,f_7,f_8,f_9\}=`$ $`\{`$ $`1.4`$, 0.55, $`1.3`$, $`13.0,0.45`$, 0.46, 0.11$`\}`$ with $`f_1`$ standing for the entire slope between 1 Ryd and $`f_4`$, whereas the point $`\{f_1,f_2,f_3,f_4,f_5,f_6,f_7\}=`$ $`\{`$ $`0.5`$, 0.25, $`1.0`$, 0.61, $`2.5`$, $`3.0,0.7`$ $`\}`$ corresponds to the AGN-type spectrum of Mathews & Ferland (1986, hereafter MF) (all energies are given in $`\mathrm{log}_{10}`$). The next step of the shape adjustment procedure is to find a direction in the factor space which leads to the UV background with more appropriate characteristics. The starting (‘null’) point is represented by the initial spectrum. A set of new trial spectra is produced by varying the factor values around the ‘null’ point in accordance with a special scheme called ‘experimental (factorial) design’. The experimental design is represented by a matrix with $`n`$ ($`n>k`$) rows containing particular values of factors (an example of the experimental design is given in Appendix). To evaluate the fitness of each trial ionizing spectrum a numerical measure $`\stackrel{~}{}`$ (usually called ‘response’) is to be defined in such a way that bigger values of $`\stackrel{~}{}`$ correspond to increasing fitness. In general, the choice of $`\stackrel{~}{}`$ occurs rather heuristically and accounts for the information obtained with the initial UV spectrum and for any a priori information (like, e.g., allowable element abundance ratios). For illustration consider the following two examples. Assume that a set of silicon (Si ii, Si iii, Si iv) and carbon (C ii, C iii, C iv) lines is observed in an absorption system, and that the initial UV spectrum overestimates the column densities of Si ii and C iv and underestimates those of Si iv and C ii. Then the response can be written as $$\stackrel{~}{}=\frac{C\text{ii}}{C\text{iv}}\times \frac{Si\text{iv}}{Si\text{ii}},$$ i.e., the search should go towards increasing the product of these ratios. Another example concerns an absorption system with the observed lines of Si iii, Si iv, and C iv. The ratio Si iii/Si iv determines the ionization parameter $`U`$. Using the corresponding ion fractions $`\mathrm{{\rm Y}}_i`$ the ratio $$\frac{\mathrm{Si}}{\mathrm{C}}=\frac{Si\text{iv}}{C\text{iv}}\frac{\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{IV}}}{\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}}}$$ can be calculated. Both observations and theoretical considerations give for this ratio a safe upper limit Si/C $`<3`$(Si/C)$``$. If, for instance, the initial UV background delivers Si/C = 10(Si/C)$``$, then the search for a new SED can be governed by the ratio $$\stackrel{~}{}=\frac{\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}}}{\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{IV}}},$$ calculated for the ionization parameter $`U`$ given by the observed column density ratio Si iii/Si iv (note that the value of $`U`$ depends on the shape of the trial spectrum). More examples are presented in the subsequent sections. The calculated responses $`\{\stackrel{~}{}_i\}_{i=1}^n`$ can be considered as points belonging to some surface in $`k`$-dimensional factor space (referred to as ‘response surface’). To determine the direction of steepest ascent the response surface should be described by some analytical function. A standard linear model is a reasonable first approximation: $$=\underset{i=1}{\overset{k}{}}\alpha _i\widehat{f}_i+\beta ,$$ (1) where $`\{\alpha _i\}_{i=1}^k`$ are factor effects, $`\widehat{f}_i`$ is the scaled and centered value of the $`i`$th factor, $`\widehat{f}_i=(f_if_{0,i})/s_i`$, $`s_i`$ is some suitable scale (range of variation) of the $`i`$th factor, and $`\beta =\alpha _0+_{i,j=1}^k\alpha _{i,j}\widehat{f}_i\widehat{f}_j`$ stands for the joint effect of the free term $`\alpha _0`$ and of the nonlinear term that contains also the interaction of factors. The coefficients $`\alpha _i`$ and $`\beta `$ and their dispersions are calculated from $`n`$ values of $``$ according to special formulas (depending on the type of experimental design used). The validity of the planar data model (1) can be checked by comparison of the estimation, $`\beta ^{}`$, with the value of $`_0`$ at the ‘null point’, i.e., when $`f_i=f_{0,i}`$. Statistical significance of the difference ($`_0\beta ^{}`$) points to non-negligible non-linear effects. In this case the higher-order data model is to be used with corresponding higher-order experimental design. Given the factor effects $`\{\alpha _i\}`$, the maximization of $``$ is obtained by moving the factor values from the ‘null’ point in the direction normal to the response surface. This movement is performed stepwise until either the desired UV background is found (i.e. the one which reproduces self-consistently all observed column densities) or the validity limit of the employed data model is reached. In the latter case the adjustment procedure should be repeated, this time with a newly obtained UV spectrum as a ‘null’ point. The initial set of factors may turn out to be redundant, i.e. some of the factor effects $`\alpha _i`$ may be statistically insignificant. This means that the corresponding features in the spectral shape do not affect the absorption line profiles included in the analysis. Such insignificant factors are removed from the data model (1) and are fixed at some appropriate level. The effects of the remaining $`m`$ factors are then recalculated using the corresponding $`m`$-factorial design. The uncertainty of the recovered spectral shape requires some comments. The spectrum of the ionizing radiation, $`J_\nu `$, defines the rates (s<sup>-1</sup>) of the photoionization processes, $`\mathrm{\Gamma }_{ij}`$, through the integral $$\mathrm{\Gamma }_{ij}=_{\nu _{ij}}^{\mathrm{}}\frac{4\pi J_\nu }{h\nu }\sigma _{ij}(\nu )𝑑\nu ,$$ (2) where $`\sigma _{ij}(\nu )`$ is the photoionization cross section at frequency $`\nu `$ for species $`i`$ in ionization state $`j`$, and $`\nu _{ij}`$ is the threshold frequency for photoionization. Thus, the SED is obtained by solving the integral equation which is a classic ill-posed problem. Its solution depends crucially from the number of constraints included in the analysis. When applied to our problem, this means that the SED is best recovered if the corresponding absorption system reveals many lines of different metals in different ionization stages. If only a few metal lines are observed, they can be used to recover the shape not in the entire range 1 Ryd $`<E<10`$ Ryd, but in some narrower regions. For example, the lines of C iii, C iv and O vi allow to estimate the SED in the vicinity of the He ii break (3 Ryd $`<E<4.5`$ Ryd). ### 2.2 The Monte Carlo Inversion procedure Absorption systems are analyzed by means of the Monte Carlo Inversion (MCI) procedure described in detail in Levshakov, Agafonova & Kegel (2000, hereafter LAK), and with modifications in Levshakov et al. (2002, 2003a,b). Here we briefly outline the basics needed to understand the results presented below in Sect. 3. The MCI is based on the assumption that all lines observed in the absorption system are formed in a continuous medium where the gas density, $`n_\mathrm{H}(x)`$, and velocity, $`v(x)`$, fluctuate from point to point giving rise to complex profiles (here $`x`$ is the space coordinate along the line of sight). The MCI also assumes that within the absorber the metal abundances are constant, the gas is optically thin for the ionizing UV radiation, and the gas is in thermal and ionization equilibrium. The intensity and the spectral shape of the background ionizing radiation are treated as external parameters. The radial velocity $`v(x)`$ and gas density $`n_\mathrm{H}(x)`$ are considered as two continuous random functions which are represented by their sampled values at equally spaced intervals $`\mathrm{\Delta }x`$. The computational procedure uses the adaptive simulated annealing. The fractional ionizations of all elements included in the analysis are computed at every space coordinate $`x`$ with the photoionization code CLOUDY. In the MCI procedure the following physical parameters are directly estimated: the mean ionization parameter $`U_0`$, the total hydrogen column density $`N_\mathrm{H}`$, the line-of-sight velocity dispersion $`\sigma _\mathrm{v}`$, and density dispersion $`\sigma _\mathrm{y}`$, of the bulk material \[$`yn_\mathrm{H}(x)/n_0`$\], and the chemical abundances $`Z_\mathrm{a}`$ of all elements involved in the analysis. With these parameters we can further calculate the column densities $`N_\mathrm{a}`$ for different species, and the mean kinetic temperature $`T_{\mathrm{kin}}`$. If the absolute intensity of the UV background is known, then the mean gas number density $`n_0`$, and the line-of-sight thickness $`L`$ of the absorber can be evaluated as well. In general, the uncertainties of the fitting parameters $`U_0`$, $`N_\mathrm{H}`$, $`\sigma _v`$, $`\sigma _y`$, and $`Z_a`$ are about 15%–20% (for data with S/N $`>\mathrm{\hspace{0.33em}30}`$) and the errors of the estimated column densities are less than 10%. However, in individual absorption systems, the accuracy of recovered values can be lower due to different reasons such as partial blending of line profiles, saturation of profiles or absence of lines of subsequent ionic transitions. The procedure of spectral shape adjustment described in the preceding section is used in a modified form. From computational point of view, it is more convenient within the MCI to evaluate responses $`\stackrel{~}{}`$ using ion fractions rather then column densities (at a given metallicity, the column density of an ion is proportional to its fraction). The procedure stops when spectral shape ensuring $`\chi ^2<\mathrm{\hspace{0.33em}1}`$ for all lines observed in the system is found. ### 2.3 Testing the SED recovering procedure In order to test the procedure of the SED adjustment we prepared a mock absorption line spectrum using the density and velocity distributions plotted in Fig. 2, the UV ionizing background shown in Fig. 3 (solid line), and the model parameters listed in Table 1, Col. 2. The resulting absorption lines after convolution of the intensities with a Gaussian-type point-spread function of FWHM = 7 km s$`^1`$ and the addition of white noise with dispersion 0.02 (S/N = 50) are shown in Fig. 4 by points with error bars. The chosen UV background can be produced by a mixture of stellar and metagalactic spectra (see, e.g., Giroux & Shull 1997). The procedure requires an input SED and as initial guess for the underlying UV background we took the metagalactic spectrum of HM shown as a long-dashed line in Fig. 3. This spectrum was defined by 7 factors $`\{`$$`f_1,f_4,f_5,f_6,f_7,f_8,f_9`$$`\}`$ with the meanings described in Sect. 2.1. The only change is for $`f_1`$ which now stands for the entire slope between 1 Ryd and the He ii break at $`f_4`$. The best MCI solution gives the physical parameters listed in Table 1, Col. 3. The corresponding synthetic absorptions are shown in Fig. 4 by the dotted lines. As expected, not all absorptions are well reproduced, since the tried UV background differs from the one used to generate the mock lines. An underestimation of the C iii and Si iii intensities along with overestimation of C iv is well seen in Fig. 4. To calculate responses, the following expression was chosen: $$\stackrel{~}{}=\mathrm{log}\left(\frac{\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{III}}}{\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{IV}}}\frac{\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{III}}}{\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{II}}}\right),$$ with an additional constraint $`\mathrm{log}(\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{III}}/\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}})>\mathrm{\hspace{0.33em}0.4}`$ \[i.e. the fractions of the corresponding ions are calculated for the value of $`U`$ determined by the condition $`\mathrm{log}(\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{III}}/\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}})>\mathrm{\hspace{0.33em}0.4}`$\]. Factors were varied at 2 levels according to a 7-factorial saturated simplex design which gave 8 new trial UV spectra (see Appendix). These spectra were inserted into CLOUDY and with the obtained fractional ionizations the responses for every spectrum were evaluated. Then the factor effects were estimated, linearity of data model checked, and a new UV background produced by moving the factor values in the direction normal to the response surface. The new solution for the background is shown in Fig. 3 by the short-dashed line, and the results of the MCI calculations are given in Table 1, Col. 4. Fitting of synthetic profiles to most absorption lines becomes much better, but C iii, C iv and Si iv lines still have large $`\chi ^2`$ values. Thus, the adjustment procedure was repeated once more, with the same set of factors, but different response $$\stackrel{~}{}=\mathrm{log}\left(\frac{\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{IV}}}{\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{II}}}\frac{\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}}}{\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{II}}}\right),$$ with constraints $`\mathrm{log}(\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{III}}/\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{IV}})>`$ 0.9 and $`\mathrm{log}(\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{III}}/\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}})>`$ 0.5. The resulting UV spectrum is shown in Fig. 3 by dots, and the corresponding synthetic profiles by the smooth lines in Fig. 4. The recovered parameters are listed in Col. 5, Table 1. Now all absorption lines are well reproduced, and, thus, the inverse procedure can be considered completed. This test demonstrates both the effectiveness of the adjustment procedure and its limitations. Starting from some standard spectral shape, we were able to recover quite correctly all significant features of the underlying ionizing spectrum such as the depth of the He ii break and its shift to the lower energies. On the other hand, some fine features of the underlying spectrum (break at 1.8 Ryd, small bump at 3 Ryd) cannot be reproduced with the adopted noise level. Higher S/N data would be more appropriate for this case. ## 3 Application to observed absorption systems Here we analyze some QSO absorption systems using the procedure described in the preceding section. All computations below were performed with laboratory wavelengths and oscillator strengths taken from Morton (2003). Solar abundances were taken from Asplund, Grevesse & Sauval (2005). ### 3.1 Absorption system at $`z_{\mathrm{abs}}`$= 2.9171 towards HE 0940–1050 This absorption system shown in Fig. 5 exhibits higher order H i Lyman series lines and transitions of C ii/C iii/C iv, and Si ii/Si iii/Si iv. The system was described in detail in Levshakov et al. (2003c) as an extremely low metallicity LLS with $`N`$(H i) $`=3\times 10^{17}`$ cm$`^2`$ and $`Z=0.001Z_{}`$. Levshakov et al. also noted that the observed line intensities were inconsistent with the HM metagalactic ionizing spectrum, and that better fitting of the observed profiles corresponded to the ionizing spectrum having a significantly enhanced He ii re-emission bump at 3 Ryd. This result was obtained by common ‘trials and errors’ method. Here we re-analyze the $`z_{\mathrm{abs}}`$= 2.9171 system using the procedure of the directed search. With the HM ionizing spectrum at $`z=3`$ as initial guess we obtained the synthetic profiles shown in Fig. 5 by dotted curves. Strong overestimation of Si ii $`\lambda 1260`$ Å is clearly seen along with underestimation of C ii $`\lambda 1334`$ Å and Si iv $`\lambda 1393`$ Å. Whether C iii $`\lambda 977`$ Å line is underestimated or not, cannot be decided unambiguously at this stage of investigation since the apparent intensity may be due to contamination by some Ly$`\alpha `$ interloper. The HM spectrum was parameterized using the same 7 factors as in the example described in Sect. 2.3. The responses were calculated in the form $$=\mathrm{log}\left(\frac{\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{II}}}{\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{IV}}}\frac{\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}}}{\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{II}}}\right).$$ The design employed was the 7-factorial saturated simplex design (see Appendix). The first-iteration ionizing spectrum revealed a sharp He ii break shifted to 3 Ryd and the He ii re-emission bump significantly increased as compared to the initial approximation (Fig. 6, dotted line). This spectrum provided better fitting to the observed profiles (e.g., all carbon lines including C iii were now described with $`\chi ^21`$), but the profiles of Si ii and Si iv remained, correspondingly, over- and underestimated. Next iteration of the SED adjustment was performed with the same set of factors, design and response augmented with constraints $`\mathrm{log}(\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{III}}/\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{IV}})<1`$ and $`\mathrm{log}(\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{III}}/\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}})>0.1`$ to enable self-consistent description of the C iii and Si iii lines. However, this step failed to produce any improvement. Attempts to include additional factors describing the He i break (see Fig. 1) or to employ other forms of the responses failed as well. The solution was found after performing model calculations (using CLOUDY) of radiation transmission through a plane-parallel absorbing cloud. For an HM-type ionizing spectrum, $`\eta `$ = $`N`$(He ii)/$`N`$(H i$`>50`$ and the absorber with $`N`$(H i) $`=3\times 10^{17}`$ cm$`^2`$, marginally optically thick in hydrogen continuum, is certainly opaque in He ii. The radiation in the He ii continuum is effectively absorbed by such cloud and a part of this radiation is reemitted as the He ii $`\lambda 304`$ Å Ly$`\alpha `$ and two-photon reemission, and He ii Balmer continuum emission. As a result, the incident HM spectrum being transmitted through the absorber with $`N`$(H i) $`=3\times 10^{17}`$ cm$`^2`$ reveals a sharp and deep break at 4 Ryd and a strong emission line at 3 Ryd of about twice the incident intensity. Thus, the increased intensity at $`E<\mathrm{\hspace{0.33em}3}`$ Ryd in the ionizing spectrum recovered from the $`z_{\mathrm{abs}}`$= 2.92 absorber is, probably, due to auto-emission of the cloud. However, in the recovered spectrum a strong intensity depression starts just above 3 Ryd, and this result has been reproduced in all trials with different responses and different sets of factors. At the same time, the model calculations show that the spectral range 3 Ryd $`<E<4`$ Ryd remains unaffected by the processes inside the cloud, i.e., if the intensity depression in this range is not present in the incident spectrum, then it does not appear in the transmitted one. This means that the depression at $`E>\mathrm{\hspace{0.33em}3}`$ Ryd in the ionizing spectrum is produced by the processes outside the cloud. We can assume that this is a manifestation of He ii Ly$`\alpha `$ $`\lambda 304`$ Å absorption arising in both smoothly distributed IGM (He ii Gunn-Peterson effect) and discrete Ly$`\alpha `$ forest clouds (line blanketing). He ii Ly$`\alpha `$ absorption in the diffuse IGM gas results in the absorption trough blueward of the resonant wavelength 304 Å (3 Ryd). This trough together with the subsequent absorption in the He ii continuum ($`E4`$ Ryd) produce a winding structure in the spectral shape at $`E>`$ 3 Ryd. As a first approximation, this structure can be described by a straight step (see Fig. 6). After this consideration, a new set of factors was defined to account for this step-like form (Fig. 6): $`f_1`$ – the slope between 1 and 3 Ryd; $`f_2`$ – the depth of the He ii absorption through; $`f_3`$ – the energy of the He ii ionization break; $`f_4`$ – the depth of the He ii ionization break; $`f_5`$ – the slope of the He ii ionization break; $`f_6,f_7`$ – the height and width of the He ii re-emission bump; $`f_8`$ – the slope after the He ii ionization break. In the following analysis, factors $`f_3`$ and $`f_8`$ turned out to have low effects and were set to 4 Ryd and –0.45, respectively. The final UV background is shown by the short-dashed line in Fig. 6. The synthetic line profiles calculated with this SED are plotted by the smooth solid lines in Fig. 5 and the corresponding physical parameters are listed in Table 2, Col. 1. The predicted column density of He ii for this absorber is $`2.8\times 10^{19}`$ cm$`^2`$ giving $`\eta `$ = 87. In the framework of our approach, the restored spectrum may be considered as some average spectrum. This spectrum reveals features stemming from both the intergalactic incident radiation (depression between 3 and 4 Ryd in Fig. 6) and local processes in the cloud itself (enhanced re-emission of He ii Ly$`\alpha `$, deep break at 4 Ryd). To reconstruct the true metagalactic spectrum, absorbers optically thin in He ii are needed, i.e. those with $`N`$(H i) $`<\mathrm{\hspace{0.33em}10}^{15}`$ cm$`^2`$. ### 3.2 Absorption system at $`z_{\mathrm{abs}}`$= 2.9659 towards Q 0347–3819 This system has been described in detail in Levshakov et al. (2003b). Since then Q 0347–3819 was re-observed at the VLT/UVES with a full wavelength coverage and longer exposure time. This allowed to obtain a high-quality spectrum with resolution of $``$ 6 km s$`^1`$and S/N = $`50100`$. Here we repeat the analysis for the new data. The hydrogen and metal absorption lines of the $`z_{\mathrm{abs}}`$= 2.9659 system are shown in Fig. 7 (histograms). Neither C ii $`\lambda 1334`$ Å (undetected), nor Si ii $`\lambda 1260`$ Å (blended) can be used. This makes the presence of C iii $`\lambda 977`$ Å line crucial for the SED estimation since Si iii, Si iv and C iv lines can be fitted with a very broad range of ionizing spectra (constrained only by the condition \[Si/C\] $`<0.5`$). Although we did not find any metal line candidate for blending with C iii, we cannot exclude blending with a Ly$`\alpha `$ forest absorption. Thus, the analysis below should be taken with caution and it is valid only if the absorption at the position of C iii is entirely due to this ion. Calculations with the HM ionizing background showed that this spectrum significantly underestimates C iii. Our trials to reconstruct the underlying continuum shape were conducted with two sets of factors, with and without accounting for the He ii Ly$`\alpha `$ absorption. In both cases it was possible to obtain the appropriate ionizing spectra. However, the spectrum estimated without He ii Ly$`\alpha `$ absorption shows significantly enhanced emission at 3 Ryd and at the same time is harder at $`E>4`$ Ryd compared to the initial HM spectrum. Since the absorber under study has $`N`$(H i) $`=3.4\times 10^{14}`$ cm$`^2`$ and is optically thin in He ii, this type of ionizing spectrum is obviously an artifact and should be rejected. The recovered spectrum with the intensity depression at 3 Ryd$`<E<`$ 4 Ryd is shown in Fig. 6 by the point-dashed line with the corresponding synthetic profiles plotted by the smooth lines in Fig. 7 (physical parameters are given in Table 2, Col. 3). Since the depth of the He ii Ly$`\alpha `$ trough (factor $`f_2`$) affects quite strongly the ion fractions, its value is estimated with accuracy better than 0.1 dex. The depth of the He ii continuum absorption (factor $`f_4`$) has an uncertainty of about 0.15 dex. The predicted He ii column density is $`2.1\times 10^{16}`$ cm$`^2`$ ($`\eta `$ = 62) and, hence, this cloud is optically thin in He ii, $`\tau _{\mathrm{He}\mathrm{II}}^{\mathrm{cont}}0.03`$. Thus, the recovered background coincides in this case with the infalling radiation spectrum and is not modified by local effects. From the depth of the He ii Ly$`\alpha `$ trough we can estimate the opacity of the intergalactic He ii Ly$`\alpha `$ absorption as $`\tau _{\mathrm{He}\mathrm{II}}`$ = $`2.1\pm 0.2`$. ### 3.3 Absorption system at $`z_{\mathrm{abs}}`$= 2.9375 towards HE 0940–1050 The absorption system at $`z_{\mathrm{abs}}`$= 2.9375 towards HE 0940–1050 reveals metal lines of C iii $`\lambda 977`$ Å, C iv $`\lambda \lambda 1548,1550`$ Å and O vi $`\lambda 1037`$ Å (O vi $`\lambda 1031`$ Å is blended, as well as N v $`\lambda \lambda 1238,1242`$ Å) (Fig. 8) . We cannot exclude that the observed O vi $`\lambda 1037`$ Å and C iii lines are contaminated by some hydrogen absorption. However, no metal line candidates for blending were found. As in the case of the $`z_{\mathrm{abs}}`$= 2.9659 system described in the preceding subsection, the analysis below is valid only if the intensities are entirely due to absorption of the corresponding ions. Two lines of one element along with one line of another can be described with a very broad range of ionizing spectra. However, different spectra will deliver different abundance ratios \[O/C\] and this can be used to constrain the allowable SEDs. Measurements of \[O/C\] in Galactic and extragalactic H ii regions and in metal poor halo stars indicate that a safe upper bound for \[O/C\] is 0.5 (Henry, Edmunds, & Köppen 2000; Akerman et al. 2004). Calculations with the HM spectrum produce for this absorption system \[O/C\] $`>1`$ suggesting that the shape of the input ionizing background could be inadequate. To meet \[O/C\] $`<0.5`$, the ionizing spectrum should be either significantly harder (by 0.5 dex) at $`E>4`$ Ryd or should have a step-like depression at $`E>`$ 3 Ryd. The second option is preferable since this step is also present in the spectrum recovered from the absorber at $`z_{\mathrm{abs}}`$= 2.9171 along the same line of sight (Sect. 3.1). The step at 3 Ryd $`<E<`$ 4 Ryd strongly affects the C iii/C iv ratio leading to the result that the observed ratio is obtained at a higher ionization parameter as compared to the HM spectrum and, hence, the corresponding O vi fraction is larger. Using the HM spectrum as an initial approximation and varying the depth of the He ii Ly$`\alpha `$ absorption trough (i.e. one-factor experiment) we adjusted the spectral shape in a way that enables the description of the observed lines with \[O/C\] $`<0.5`$. The obtained spectral shape is shown in Fig. 6 (long-dashed line), with the synthetic profiles plotted by the smooth lines in Fig. 7 and the physical parameters given in Table 2, Col. 4. The uncertainty in the step depth is of 0.1 dex (given the energy above 4 Ryd at the level of the HM spectrum). The predicted column density of the once ionized helium $`N`$(He ii) $`=8.2\times 10^{16}`$ cm$`^2`$ translates to $`\eta `$ = 146 which is almost 2.5 times higher as in the $`z_{\mathrm{abs}}`$= 2.9659 system. ### 3.4 Absorption system at $`z_{\mathrm{abs}}`$= 1.9426 towards J 2233–606 This system has been described by Prochaska & Burles (1999), D’Odorico & Petitjean (2001) and Levshakov et al. (2002). A large wavelength coverage provided by combining the VLT/UVES and HST/STIS spectra, allows to identify many higher order hydrogen lines and numerous metal transitions (Fig. 9). Unfortunately, the important C ii $`\lambda 1334`$ Å line is blended with a Ly-$`\alpha `$ forest absorption, and C ii $`\lambda 1036`$ Å is very noisy, but this system is nevertheless worth being analyzed due to the rare occasion of simultaneous presence of low- (Si ii, Mg ii) and high ionization lines like N v. Calculations with the $`z=2`$ HM ionizing spectrum (Fig. 10, solid line) describe most of the observed line profiles except N iii $`\lambda 989`$ Å and N v $`\lambda \lambda 1238,1242`$ Å which come over- and underestimated, respectively. The SED was adjusted in two experiments with factor sets both accounting for the Gunn-Peterson (GP) absorption and without it, and with the response maximizing the product of the ratios $`\mathrm{{\rm Y}}_{\mathrm{N}\mathrm{V}}/\mathrm{{\rm Y}}_{\mathrm{N}\mathrm{III}},\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{II}}/\mathrm{{\rm Y}}_{\mathrm{C}\mathrm{IV}}`$, and $`\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{IV}}/\mathrm{{\rm Y}}_{\mathrm{Si}\mathrm{III}}`$. All trials showed that to reproduce the nitrogen lines the ionizing background should be significantly harder at $`E>4`$ Ryd than the initial HM spectrum. Unfortunately, because of a low quality of the C ii $`\lambda \lambda 1334,1036`$ Å lines, it turned out to be impossible to distinguish between the ionizing spectra either with or without the GP absorption. In fact, both recovered SEDs (shown by the dotted and dashed lines in Fig. 10) give identical fitting to the data. The synthetic profiles are shown by the smooth lines in Fig. 9. However, an upper limit can be set for a putative GP He ii absorption at $`z=2`$: $`\tau _{\mathrm{He}\mathrm{II}}<1.8`$. The predicted column density of He ii of $`1.1\times 10^{18}`$ cm$`^2`$ ($`\eta `$ = 44) gives the optical depth in the He ii continuum of 1.5. In principle, taken at face value such optical depth can soften noticeably the incident ionizing continuum at $`E>`$ 4 Ryd. However, the observed lines of N iii and N v clearly require a hard spectrum in this energy range. This contradiction can be explained in two ways: either the N v lines arise in the external regions of the absorbing cloud where the incident ionizing continuum is not yet distorted by the He ii continuum absorption, or the density variations along the line of sight make the effective optical depth smaller. In both cases we can conclude that the metagalactic ionizing spectrum at $`z2`$ is harder at $`E>4`$ Ryd as compared to the spectrum predicted by HM. ## 4 Conclusions We have proposed a method to estimate the shape of the ionizing continuum from metal lines observed in the intervening absorption systems. The implementation includes the following steps: 1. parameterization of the spectral shape by means of a set of factors based on physical processes relevant for the IGM; 2. choice of some quantitative measure (response) to estimate the fitness of a trial shape; 3. generation of trial shapes by randomization of the factor values according to some experimental design; 4. evaluation of the factor effects; 5. estimation of the optimal shape by moving in factor space towards the ascending fitness. The result depends on the number of metal lines involved in the analysis: in general, the more lines of different ionic transitions of different elements are detected in an absorption system the higher is the accuracy of the recovered spectral shape. In some cases additional information concerning, for instance, element abundance ratios can significantly tighten the factor values, in particular when only a few metal lines are available. Although the main objective of this work is to illustrate how the proposed approach can be used in practice, some physical results are worth mentioning as well. They are as follows: 1. The metagalactic ionizing spectrum at redshift $`z3`$ has breaks at 3 and 4 Ryd and is quite hard at $`E>4`$ Ryd, at least as hard as a model metagalactic spectrum of Haardt & Madau (1996). 2. The intensity decrease between 3 and 4 Ryd is probably produced mostly by He ii Ly$`\alpha `$ absorption in the intergalactic diffuse gas since line-blanketing from discrete Ly$`\alpha `$ forest clouds cannot account for significant He ii opacity with such hard ionizing spectrum (Fardal, Giroux & Shull 1998; Zheng, Davidsen & Kriss 1998). Thus, the intensity depression at 3 Ryd $`<E<4`$ Ryd in the spectrum of the metagalactic ionizing radiation may be an imprint of a true He ii Gunn-Peterson effect. 3. The He ii opacity estimated from the depth of the He ii Ly$`\alpha `$ absorption trough is $`\tau _{\mathrm{He}\mathrm{II}}`$ $``$ 2. This should be considered as an average optical depth since in our approach the He ii Ly$`\alpha `$ absorption is approximated by a straight step. As known from observations, He ii opacity at $`z3`$ fluctuates revealing both absorption troughs and opacity gaps. For reference, $`\tau _{\mathrm{He}\mathrm{II}}`$ $``$ 4 is measured in absorption troughs at $`2.77<z<2.87`$ in the spectrum of HE 2347–4342 (Zheng et al. 2004) along with $`\tau _{\mathrm{He}\mathrm{II}}<0.5`$ in the opacity gaps at $`z_{\mathrm{abs}}`$= 2.817 and 2.866 (Reimers et al. 1997). 4. A putative galactic contribution to the ionizing background spectrum (e.g. Pettini et al. 2001; Ciardi, Bianchi & Ferrara 2002; Fujita et al. 2003), if present, would significantly soften the spectrum at $`E>4`$ Ryd. We do not find any traces of soft component in the recovered spectra: even in the absorption systems at $`z=2.96`$ and $`z=1.94`$, where high metallicity supposes their kinship with galaxies, the observed line intensities point to a hard ionizing background. Thus, we can conclude that the ionizing background at $`z3`$ and $`z2`$ is dominated by QSOs. ###### Acknowledgements. The work of I.I.A. and S.A.L. is supported by the RFBR grant 03-02-17522 and by the RLSS grant 1115.2003.2. ## Appendix A An example of experimental design for 7 factors The factors are tested at two levels: $`f_i=f_{0,i}\pm \sigma _i`$, where $`\sigma _i`$ is a suitable scale. Usually $`\sigma _i(0.10.2)f_{0,i}`$, in order to maintain the linear dependence of the response $`\stackrel{~}{}`$ on the factors. The normalized and centered factor values at the upper and lower levels are, respectively, $`+1`$ and $`1`$. Then a 7-factorial saturated simplex design can be defined by the following symmetric and orthogonal matrix $`\{\widehat{f}_{ij}\}`$ (Nalimov 1971): \| \| $`\widehat{f}_1`$ \| $`\widehat{f}_2`$ \| $`\widehat{f}_3`$ \| $`\widehat{f}_4`$ \| $`\widehat{f}_5`$ \| $`\widehat{f}_6`$ \| $`\widehat{f}_7`$ \| $`\stackrel{~}{}`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| +1 \| –1 \| –1 \| –1 \| +1 \| +1 \| +1 \| –1 \| $`_1`$ \| \| +1 \| +1 \| –1 \| –1 \| –1 \| –1 \| +1 \| +1 \| $`_2`$ \| \| +1 \| –1 \| +1 \| –1 \| –1 \| +1 \| –1 \| +1 \| $`_3`$ \| \| +1 \| +1 \| +1 \| –1 \| +1 \| –1 \| –1 \| –1 \| $`_4`$ \| \| +1 \| –1 \| –1 \| +1 \| +1 \| –1 \| –1 \| +1 \| $`_5`$ \| \| +1 \| +1 \| –1 \| +1 \| –1 \| +1 \| –1 \| –1 \| $`_6`$ \| \| +1 \| –1 \| +1 \| +1 \| –1 \| –1 \| +1 \| –1 \| $`_7`$ \| \| +1 \| +1 \| +1 \| +1 \| +1 \| +1 \| +1 \| +1 \| $`_8`$ \| The first column is added to estimate the free term $`\beta `$ in the model (1). Such experimental design is called a ‘two level saturated plan’ since here all degrees of freedom $`N=8`$ are used to estimate 8 regression coefficient ($`\{\alpha _i\}_{i=1}^7`$ and $`\beta `$): $$\alpha _i=\frac{1}{N}\underset{j=1}{\overset{N}{}}\widehat{f}_{ji}_j,$$ and $$\beta =\frac{1}{N}\underset{j=1}{\overset{N}{}}_j.$$ Other plans with different numbers of levels and treatments (rows in the matrix) can be used as well. After the regression coefficients for the model (1) are found, new factor values can be calculated from $$\widehat{f}_i^{\mathrm{new}}=f_{0,i}+\mathrm{\Delta }\sigma _i\mu _i,$$ where $`\mu _i=\alpha _i/\sqrt{_{j=1}^7\alpha _j^2}`$ are direction cosines and $`\mathrm{\Delta }`$ is the increment of the response. The value of $`\mathrm{\Delta }`$ is restricted by the condition $`\chi ^21`$ for every absorption line (or its portion) involved in the optimization procedure.
warning/0506/hep-th0506063.html	ar5iv	text	# 1 Introduction ## 1 Introduction The idea for the correspondence between the large N limit of gauge theories and string theory was proposed over thirty years ago but its realization was given when Maldacena conjectured the AdS/CFT correspondence . Since then this became a major research area and many fascinating discoveries were made in the last years. These include the discovery of Gubser, Klebanov and Polyakov that the energy of certain string configurations in the limit of large quantum numbers reproduces the behavior of the anomalous dimension of the corresponding SYM operator. Soon after that Minahan and Zarembo provided a way to compute the anomalous dimension of a certain dilatation operator by presenting it as the Hamiltonian of an integrable spin chain. These two discoveries opened the door towards many qualitative and quantative checks of the AdS/CFT beyond the supergravity approximation. On the string side of the correspondence many semiclassical string configurations on $`AdS_5\times S^5`$ were studied and the classical energies as well as their quantum corrections were obtained <sup>2</sup><sup>2</sup>2see and for a review -. On the gauge theory side the anomalous dimensions of the dilatation operator in certain sectors of the $`𝒩=4`$ SYM were found (see and references therein). Moreover these dimensions coincide with the energies of the corresponding semiclassical strings, thus providing remarkable quantative proofs of the string theory/gauge theory correspondence -. Nevertheless the $`AdS_5\times S^5`$/$`𝒩=4`$ SYM is the main and best explored example. Since the $`𝒩=4`$ superconformal gauge theory is not appropriate phenomologically we would like to extend the above ideas to less supersymmetric Yang-Mills theories. Such attempts were made by studying semiclassical strings in less supersymmetric backgrounds (Maldacena-Nunez, Pilch-Warner and other confining and warped geometries -) but it was not completely clear how to reproduce the string theory results from the SYM side. Fortunately such a possibility emerged when Lunin and Maldacena found the gravity background dual to the Leigh-Strassler $`\beta `$-deformation of $`𝒩=4`$ SYM. A quantative check of this correspondence for the $`su(2)`$ sector was made in , thus giving the hope that we can extend the remarkable previous results to this case. Moreover in the integrability of the string Hamiltonian on the Lunin-Maldacena background was proven by finding a Lax pair. This suggests that the interplay between integrable structures in $`AdS_5\times S^5`$ and $`𝒩=4`$ SYM is also present in this less supersymmetric case. These recent developments give us the motivation to investigate semiclassical strings on the Lunin-Maldacena background. In a simple two spin string ansatz was considered. In this paper we would like to investigate multispin rotating string solutions, having angular momenta both in the $`AdS_5`$ and the $`S^5`$ part of the background. Since the anti-de Sitter piece of the metric stays undeformed under the TsT transformation generating the Lunin-Maldacena background we expect that the string motion there will not differ from the $`AdS_5\times S^5`$ case. However the case of rotating string with three spins in the deformed $`S^5`$ part will lead to some non-trivial results for the energy. We will consider the following cases: i)three spins on the deformed $`S^5`$ ($`J_1,J_2,J_3`$); ii)two spins on $`AdS_5`$ and two spins on the deformed $`S^5`$ ($`S_1,S_2,J_1,J_2`$) and iii)the most general configuration of five spins ($`S_1,S_2,J_1,J_2,J_3`$). In all of these cases we compute the energy of the rotating string in terms of the angular momenta and the string winding numbers. If we take the limit $`\stackrel{~}{\gamma }0`$ our results reproduce those found for the undeformed $`AdS_5\times S^5`$ case . This should be expected since this is exactly the limit in which the Lunin-Maldacena background reduces to the usual $`AdS_5\times S^5`$. We would also like to note that it should be possible to reproduce our results from the SYM side. For example the three $`S^5`$ spin solution should be dual to operators from the $`su(3)`$ sector, their anomalous dimensions can be found by a spin chain computation in the same spirit as this was done for the $`su(2)`$ case in <sup>3</sup><sup>3</sup>3After this paper was completed an interesting paper treating the three spin sector in the case of completely broken supersymmetry appeared .. In the next section we will briefly review the form of the Lunin-Maldacena background and find the energy for the rotating three spin string ansatz. In section 3 the case of two spins on $`AdS_5`$ and two spins on $`S^5`$ will be presented. In section 4 we will find the general rotating string solution with five spins and compute its energy. Finally in the last section we will present our conclusions and some open problems. ## 2 Three spin string solution in Lunin-Maldacena background We will begin by presenting the form of the supergravity background found by Lunin and Maldacena: $$ds_{str}^2=R^2\sqrt{H}\{ds_{AdS_5}^2+\underset{i=1}{\overset{3}{}}(d\rho _i^2+G\rho _i^2d\varphi _i^2)+(\stackrel{~}{\gamma }^2+\stackrel{~}{\sigma }^2)G\rho _1^2\rho _2^2\rho _3^2(\underset{i=1}{\overset{3}{}}d\varphi _i)^2\}$$ (2.1) where $$\begin{array}{c}\frac{1}{G}=1+(\stackrel{~}{\gamma }^2+\stackrel{~}{\sigma }^2)(\rho _1^2\rho _2^2+\rho _1^2\rho _3^2+\rho _2^2\rho _3^2)H=1+\stackrel{~}{\sigma }^2(\rho _1^2\rho _2^2+\rho _1^2\rho _3^2+\rho _2^2\rho _3^2)\\ \\ B_2=R^2(\gamma Gw_2\sigma w_1d\psi )\psi =\frac{\varphi _1+\varphi _2+\varphi _3}{3}\\ \\ dw_1=\mathrm{cos}\alpha \mathrm{sin}^3\alpha \mathrm{sin}\theta \mathrm{cos}\theta d\alpha d\theta w_2=\rho _1^2\rho _2^2d\varphi _1d\varphi _2+\rho _2^2\rho _3^2d\varphi _2d\varphi _3+\rho _3^2\rho _1^2d\varphi _3d\varphi _1\\ \\ \rho _1=\mathrm{sin}\alpha \mathrm{cos}\theta \rho _2=\mathrm{sin}\alpha \mathrm{sin}\theta \rho _3=\mathrm{cos}\alpha \end{array}$$ (2.2) Now we are going to write the Polyakov action for strings staying at the center of $`AdS_5`$ and moving on the deformed five sphere (i.e. strings on $`R_t\times S_\beta ^5`$).We will also suppose that the deformation parameter is real ($`\stackrel{~}{\gamma }0`$, $`\stackrel{~}{\sigma }=0`$). $$\begin{array}{c}S=\frac{R^2}{2}\frac{d\tau d\sigma }{2\pi }[\gamma ^{\alpha \beta }\sqrt{H}(_\alpha t_\beta t+_\alpha \alpha _\beta \alpha +\mathrm{sin}^2\alpha _\alpha \theta _\beta \theta \hfill \\ \\ \hfill +G\mathrm{sin}^2\alpha \mathrm{cos}^2\theta _\alpha \varphi _1_\beta \varphi _1+G\mathrm{sin}^2\alpha \mathrm{sin}^2\theta _\alpha \varphi _2_\beta \varphi _2+G\mathrm{cos}^2\alpha _\alpha \varphi _3_\beta \varphi _3\\ \\ \hfill +\stackrel{~}{\gamma }^2G\mathrm{sin}^4\alpha \mathrm{cos}^2\alpha \mathrm{sin}^2\theta \mathrm{cos}^2\theta (_\alpha \varphi _1+_\alpha \varphi _2+_\alpha \varphi _3)(_\beta \varphi _1+_\beta \varphi _2+_\beta \varphi _3))\\ \\ \hfill 2\stackrel{~}{\gamma }Gϵ^{\alpha \beta }(\mathrm{sin}^4\alpha \mathrm{sin}^2\theta \mathrm{cos}^2\theta _\alpha \varphi _1_\beta \varphi _2+\mathrm{sin}^2\alpha \mathrm{cos}^2\alpha \mathrm{sin}^2\theta _\alpha \varphi _2_\beta \varphi _3\\ \\ \hfill +\mathrm{sin}^2\alpha \mathrm{cos}^2\alpha \mathrm{cos}^2\theta _\alpha \varphi _3_\beta \varphi _1)]\end{array}$$ (2.3) Let us consider the following rotating string ansatz: $$\varphi _1=\omega _1\tau +m_1\sigma \varphi _2=\omega _2\tau +m_2\sigma \varphi _3=\omega _3\tau +m_3\sigma $$ (2.4) We impose also a constant radii condition, namely $`\alpha =\theta ={\displaystyle \frac{\pi }{4}}`$ and the global time is expressed through the world sheet time as $`t=\kappa \tau `$. Then the equations of motion for $`\alpha `$ and $`\theta `$ are simply the following relations between the frequencies $`\omega _i`$, the winding numbers $`m_i`$ and the real deformation parameter $`\stackrel{~}{\gamma }`$: $$\omega _1^2m_1^2\omega _2^2+m_2^2+\frac{\stackrel{~}{\gamma }}{2}(\omega _3m_2\omega _2m_3\omega _3m_1+\omega _1m_3)=0$$ (2.5) and $$\begin{array}{c}m_1^2\omega _1^2\omega _2^2+m_2^2+2\omega _3^22m_3^2\frac{\stackrel{~}{\gamma }^2}{8}(\omega _1+\omega _2+\omega _3)^2+\frac{\stackrel{~}{\gamma }^2}{8}(m_1+m_2+m_3)^2\hfill \\ \\ \hfill +\stackrel{~}{\gamma }(\omega _2m_1\omega _1m_2)=0\end{array}$$ (2.6) We should also explore the Virasoro constraints which in our case have the following form: $$\omega _1m_1+\omega _2m_2+2\omega _3m_3+\frac{\stackrel{~}{\gamma }^2}{8}(\omega _1m_1+\omega _2m_2+\omega _3m_3+\omega _1m_2+\omega _1m_3+\omega _2m_3)=0$$ (2.7) and $$\begin{array}{c}\frac{4\kappa ^2}{G}=\omega _1^2+\omega _2^2+2\omega _3^2+m_1^2+m_2^2+2m_3^2\hfill \\ \\ \hfill +\frac{\stackrel{~}{\gamma }^2}{8}(\omega _1^2+\omega _2^2+\omega _3^2+2\omega _1\omega _2+2\omega _1\omega _3+2\omega _2\omega _3+m_1^2+m_2^2+m_3^2+2m_1m_2+2m_1m_3+2m_2m_3)\end{array}$$ (2.8) It is easily seen that if $$m_1=m_2=m_3=m\omega _1=\omega _2=\omega _3=\omega $$ (2.9) (2.5) and (2.7) are satisfied, from (2.6) follows that $`3\omega =m`$ and the second Virasoro constraint (2.8) gives and expression for $`\kappa `$ which is related to the energy in the following way (we just note that $`R^2=\sqrt{\lambda }`$): $$E=\sqrt{\lambda }=\frac{R^2}{2\pi }_0^{2\pi }𝑑\sigma \kappa =R^2\kappa $$ (2.10) Before we calculate the energy let us first compute the three conserved charges corresponding to the three angle variables in our problem. $$\begin{array}{c}𝒥_1=\frac{J_1}{\sqrt{\lambda }}=\frac{G}{2}(\omega +\frac{3\stackrel{~}{\gamma }^2}{8}\omega +\frac{3\stackrel{~}{\gamma }}{4}m)\\ \\ 𝒥_2=\frac{J_2}{\sqrt{\lambda }}=\frac{G}{2}(\omega +\frac{3\stackrel{~}{\gamma }^2}{8}\omega \frac{3\stackrel{~}{\gamma }}{4}m)\\ \\ 𝒥_3=\frac{J_3}{\sqrt{\lambda }}=G(\omega +\frac{3\stackrel{~}{\gamma }^2}{16}\omega )\end{array}$$ (2.11) And thus for the full spin of our system we find: $$𝒥=𝒥_1+𝒥_2+𝒥_3=2G\omega +\frac{9\stackrel{~}{\gamma }^2}{16}G\omega =2G\omega +\frac{3\stackrel{~}{\gamma }^2}{16}Gm$$ (2.12) The equation for $`\kappa `$ is simply: $$\kappa =\sqrt{G(\omega ^2+m^2+\frac{9\stackrel{~}{\gamma }^2}{32}\omega ^2+\frac{\stackrel{~}{\gamma }^2}{32}m^2)}$$ (2.13) Thus for the energy we end up with the following, messy at first sight,expression: $$\begin{array}{c}E=\sqrt{\lambda }[(1+\frac{5\stackrel{~}{\gamma }^2}{16})\frac{𝒥^2}{4}\frac{3\stackrel{~}{\gamma }^2}{32}(1+\frac{9\stackrel{~}{\gamma }^2}{16})m𝒥\hfill \\ \\ \hfill +(1+\frac{\stackrel{~}{\gamma }^2}{32}+\frac{9\stackrel{~}{\gamma }^4}{64(16+5\stackrel{~}{\gamma }^2)}+\frac{81\stackrel{~}{\gamma }^6}{2056(16+5\stackrel{~}{\gamma }^2)})m^2]^{1/2}\end{array}$$ (2.14) Although it looks complicated this expression reproduces the result for the energy of a three spin string in pure $`AdS_5\times S^5`$ if we take the limit $`\stackrel{~}{\gamma }0`$ (as should be expected), namely: $$E=\sqrt{\lambda }\sqrt{\frac{𝒥^2}{4}+m^2}$$ (2.15) ## 3 Two spins in both $`AdS_5`$ and $`S_{\stackrel{~}{\gamma }}^5`$ Here we will investigate semiclassical strings with two spins on the $`AdS_5`$ part and two spins on the deformed $`S^5`$ part of the Lunin-Maldacena background. The relevant Polyakov action is: $$\begin{array}{c}S=\frac{R^2}{2}\frac{d\tau d\sigma }{2\pi }[\gamma ^{\alpha \beta }\sqrt{H}(\mathrm{cosh}^2\rho _\alpha t_\beta t+_\alpha \rho _\beta \rho +\mathrm{sinh}^2\rho _\alpha \psi _\beta \psi \hfill \\ \\ \hfill +\mathrm{sinh}^2\rho \mathrm{cos}^2\psi _\alpha \psi _1_\beta \psi _1+\mathrm{sinh}^2\rho \mathrm{sin}^2\psi _\alpha \psi _2_\beta \psi _2+_\alpha \theta _\beta \theta \\ \\ \hfill +G\mathrm{cos}^2\theta _\alpha \varphi _1_\beta \varphi _1+G\mathrm{sin}^2\theta _\alpha \varphi _2_\beta \varphi _2)2\stackrel{~}{\gamma }Gϵ^{\alpha \beta }\mathrm{sin}^2\theta \mathrm{cos}^2\theta _\alpha \varphi _1_\beta \varphi _2]\end{array}$$ (3.1) We have imposed $`\alpha =\pi /2`$ and thus: $$G^1=1+\frac{1}{4}(\stackrel{~}{\gamma }^2+\stackrel{~}{\sigma }^2)\mathrm{sin}^22\theta H=1+\frac{\stackrel{~}{\sigma }^2}{4}\mathrm{sin}^22\theta $$ (3.2) Now we will assume that $`\stackrel{~}{\sigma }=0`$, i.e. we will work in the real deformed $`AdS_5\times S^5`$ background. It easily checked that the following ansatz is compatible with the string equations of motion: $$\begin{array}{c}t=\kappa \tau \rho =const\psi =\frac{\pi }{4}\psi _1=\nu _1\tau +n_1\sigma \psi _2=\nu _2\tau +n_2\sigma \\ \\ \alpha =\frac{\pi }{2}\theta =\frac{\pi }{4}\varphi _1=\omega _1\tau +m_1\sigma \varphi _2=\omega _2\tau +m_2\sigma \varphi _3=0\end{array}$$ (3.3) From the equations of motion for $`\rho `$, $`\psi `$ and $`\theta `$ follow some relations between the winding numbers and the frequencies: $$\omega _1^2m_1^2=\omega _2^2m_2^2\nu _1^2n_1^2=\nu _2^2n_2^2\kappa ^2=\nu _1^2n_1^2$$ (3.4) The Virasoro constraints of our system adopt the following form: $$\begin{array}{c}\frac{\mathrm{sinh}^2\rho }{2}(\nu _1n_1+\nu _2n_2)+\frac{G}{2}(\omega _1m_1+\omega _2m_2)=0\\ \\ \kappa ^2\mathrm{cosh}^2\rho =\frac{\mathrm{sinh}^2\rho }{2}(\nu _1^2+\nu _2^2+n_1^2+n_2^2)+\frac{G}{2}(\omega _1^2+\omega _2^2+m_1^2+m_2^2)\end{array}$$ (3.5) The equations of motion and the first Virasoro constraint are satisfied if we choose: $$\nu =\nu _1=\nu _2\omega =\omega _1=\omega _2n=n_1=n_2m=m_1=m_2$$ (3.6) For the two $`AdS_5`$ angular momenta ($`𝒮_1`$, $`𝒮_2`$) and the two $`S^5`$ angular momenta ($`𝒥_1`$, $`𝒥_2`$) we obtain: $$𝒮_1=𝒮_2=\frac{\mathrm{sinh}^2\rho }{2}\nu 𝒥_1=𝒥_2=\frac{G}{2}\omega \frac{\stackrel{~}{\gamma }G}{4}m$$ (3.7) The full $`AdS_5`$ and $`S^5`$ angular momenta are simply: $$𝒮=𝒮_1+𝒮_2=\mathrm{sinh}^2\rho \nu 𝒥=𝒥_1+𝒥_2=G\omega \frac{\stackrel{~}{\gamma }G}{2}m$$ (3.8) This leads to the following relation: $$\omega =𝒥+\frac{\stackrel{~}{\gamma }}{2}(m+\frac{\stackrel{~}{\gamma }}{2}𝒥)$$ (3.9) And thus the second Virasoro constraint can be expressed as: $$\kappa ^2\mathrm{cosh}^2\rho =\mathrm{sinh}^2\rho (\nu ^2+n^2)+𝒥^2+(m+\frac{\stackrel{~}{\gamma }}{2}𝒥)^2$$ (3.10) The full energy of our system is $`=\kappa \mathrm{cosh}^2\rho `$ and we find the following relation: $$\frac{}{\kappa }\frac{𝒮}{\nu }=1$$ (3.11) Using this relation, the second Virasoro constraint and the relation coming from the equation of motion for $`\rho `$ we end up with: $$2\kappa \kappa ^2=2\sqrt{n^2+\kappa ^2}𝒮+𝒥^2+(m+\frac{\stackrel{~}{\gamma }}{2}𝒥)^2$$ (3.12) This expression reduces to the analogous expression found by Arutyunov, Russo and Tseytlin if we take the limit $`\stackrel{~}{\gamma }0`$. This is exactly what one should expect because this limit reproduces the well known $`AdS_5\times S^5`$ background, which is the case considered in . ## 4 Generalized rotating string Considering the results from the preceding sections it seems natural to combine them in order to investigate the most general rotating string ansatz. This is the case of two spins in the $`AdS_5`$ part and three spins in the $`S^5`$ part of the geometry ($`S_1,S_2,J_1,J_2,J_3`$). Following the same procedure we start with the relevant string action: $$\begin{array}{c}S=\frac{R^2}{2}\frac{d\tau d\sigma }{2\pi }[\gamma ^{\alpha \beta }\sqrt{H}(\mathrm{cosh}^2\rho _\alpha t_\beta t+_\alpha \rho _\beta \rho +\mathrm{sinh}^2\rho _\alpha \psi _\beta \psi \hfill \\ \\ \hfill +\mathrm{sinh}^2\rho \mathrm{cos}^2\psi _\alpha \psi _1_\beta \psi _1+\mathrm{sinh}^2\rho \mathrm{sin}^2\psi _\alpha \psi _2_\beta \psi _2+_\alpha \alpha _\beta \alpha +\mathrm{sin}^2\alpha _\alpha \theta _\beta \theta \\ \\ \hfill +G\mathrm{sin}^2\alpha \mathrm{cos}^2\theta _\alpha \varphi _1_\beta \varphi _1+G\mathrm{sin}^2\alpha \mathrm{sin}^2\theta _\alpha \varphi _2_\beta \varphi _2+G\mathrm{cos}^2\alpha _\alpha \varphi _3_\beta \varphi _3\\ \\ \hfill +\stackrel{~}{\gamma }^2G\mathrm{sin}^4\alpha \mathrm{cos}^2\alpha \mathrm{sin}^2\theta \mathrm{cos}^2\theta (_\alpha \varphi _1+_\alpha \varphi _2+_\alpha \varphi _3)(_\beta \varphi _1+_\beta \varphi _2+_\beta \varphi _3))\\ \\ \hfill 2\stackrel{~}{\gamma }Gϵ^{\alpha \beta }(\mathrm{sin}^4\alpha \mathrm{sin}^2\theta \mathrm{cos}^2\theta _\alpha \varphi _1_\beta \varphi _2+\mathrm{sin}^2\alpha \mathrm{cos}^2\alpha \mathrm{sin}^2\theta _\alpha \varphi _2_\beta \varphi _3\\ \\ \hfill +\mathrm{sin}^2\alpha \mathrm{cos}^2\alpha \mathrm{cos}^2\theta _\alpha \varphi _3_\beta \varphi _1)]\end{array}$$ (4.1) We will again examine the case of real deformation parameter and use the following ansatz: $$\begin{array}{c}t=\kappa \tau \rho =const\psi =\frac{\pi }{4}\psi _1=\nu _1\tau +n_1\sigma \psi _2=\nu _2\tau +n_2\sigma \\ \\ \alpha =\frac{\pi }{4}\theta =\frac{\pi }{4}\varphi _1=\omega _1\tau +m_1\sigma \varphi _2=\omega _2\tau +m_2\sigma \varphi _3=\omega _3\tau +m_3\sigma \end{array}$$ (4.2) From the equations of motion for $`\rho `$, $`\psi `$, $`\alpha `$ and $`\theta `$ we can extract the following relations between the frequencies and the winding numbers. $$\begin{array}{c}\nu _1^2n_1^2=\nu _2^2n_2^2\kappa ^2=\nu _1^2n_1^2\\ \\ \omega _1^2m_1^2\omega _2^2+m_2^2+\frac{\stackrel{~}{\gamma }}{2}(\omega _3m_2\omega _2m_3\omega _3m_1+\omega _1m_3)=0\\ \\ m_1^2\omega _1^2\omega _2^2+m_2^2+2\omega _3^22m_3^2\frac{\stackrel{~}{\gamma }^2}{8}(\omega _1+\omega _2+\omega _3)^2\\ \\ +\frac{\stackrel{~}{\gamma }^2}{8}(m_1+m_2+m_3)^2+\stackrel{~}{\gamma }(\omega _2m_1\omega _1m_2)=0\end{array}$$ (4.3) We should also impose the Virasoro constraints, they will provide one more relation between the parameters and an equation for $`\kappa `$. We will retain from presenting their explicit form here but he expression are analogous to the previous two cases. In order to satisfy the constraints and the equations of motion we should choose: $$\nu =\nu _1=\nu _2\omega =\omega _1=\omega _2=\omega _3n=n_1=n_2m=m_1=m_2=m_3$$ (4.4) It is again straightforward to compute the angular momenta in both parts of the background: $$\begin{array}{c}𝒥=𝒥_1+𝒥_2+𝒥_3=2G\omega +\frac{9\stackrel{~}{\gamma }^2}{16}G\omega =2G\omega +\frac{3\stackrel{~}{\gamma }^2}{16}Gm\\ \\ 𝒮=𝒮_1+𝒮_2=\mathrm{sinh}^2\rho \nu \end{array}$$ (4.5) We just remind that in this case $`G^1=1+\frac{5}{16}\stackrel{~}{\gamma }^2`$. The equation for $`\kappa `$ following from the second Virasoro constraint is: $$\kappa ^2\mathrm{cosh}^2\rho =\mathrm{sinh}^2\rho (\nu ^2+n^2)+G(\omega ^2+m^2+\frac{9\stackrel{~}{\gamma }^2}{32}\omega ^2+\frac{\stackrel{~}{\gamma }^2}{32}m^2)$$ (4.6) The energy of the rotating string is $`=\kappa \mathrm{cosh}^2\rho `$ and it is related to the $`AdS`$ angular momentum as: $$\frac{}{\kappa }\frac{𝒮}{\nu }=1$$ (4.7) And we can extract an analogous to (3.12) expression: $$\begin{array}{c}2\kappa \kappa ^2=2\sqrt{n^2+\kappa ^2}𝒮+\left(1+\frac{5\stackrel{~}{\gamma }^2}{16}\right)\frac{𝒥^2}{4}\frac{3\stackrel{~}{\gamma }^2}{32}\left(1+\frac{9\stackrel{~}{\gamma }^2}{16}\right)m𝒥\hfill \\ \\ \hfill +\left(1+\frac{\stackrel{~}{\gamma }^2}{32}+\frac{9\stackrel{~}{\gamma }^4}{64(16+5\stackrel{~}{\gamma }^2)}+\frac{81\stackrel{~}{\gamma }^6}{2056(16+5\stackrel{~}{\gamma }^2)}\right)m^2\end{array}$$ (4.8) This expression again reduces to the one known from the undeformed $`AdS_5\times S^5`$ case if we take the limit $`\stackrel{~}{\gamma }0`$. ## 5 Conclusions and Outlook In this paper we have studied rotating strings configurations in the recently found Lunin-Maldacena background. These semiclassical strings are an important tool for proving the AdS/CFT beyond the supergravity approximation. We have found the energy of different rotating strings in terms of the angular momenta and the string winding numbers. In the limit of zero deformation parameter we reproduce the well known results from the $`AdS_5\times S^5`$ case, this should be expected because this is the limit in which the Lunin-Maldacena background reduces to the usual $`AdS_5\times S^5`$. Our work can be extended in several ways <sup>4</sup><sup>4</sup>4For other related work on the subject see - . First of all we can work with the most general ansatz with different frequencies and winding numbers and find the energy behavior in terms of the angular momenta. It is also important to consider fluctuations around this classical solutions. They will provide the corrections to the above found energies and also clarify the stability of these solutions. Maybe the most important open problem is to reproduce our results from the gauge theory side. We believe that this could be done on the level of effective actions . This was the approach in , where an exact agreement was found for the $`su(2)`$ sector and we expect that this agreement could be extended for the $`su(3)`$ sector considered in the second section. There is one more interesting class of semiclassical strings - pulsating strings. It is worth investigating such pulsating solutions in the Lunin-Maldacena background and see how the deformation affects the form of the solution. We plan to address this issue in a future paper. Acknowledgments We would like to thank Carolos Nunez for useful suggestions and comments. The work of N.P.B. was supported by an EVRIKA foundation educational award.
warning/0506/astro-ph0506686.html	ar5iv	text	# On the nature of nearby GRB/SN host galaxies ## 1 Introduction It is now well established that long-duration $`\gamma `$-ray bursts (GRBs) coincide with the explosions of certain massive stars, i.e., with a subset of very energetic core-collapse supernovae (SNe). In particular, three such cases must be regarded as secure. The first was the association of SN 1998bw with GRB 980425 (Galama et al., 1998). This connection was, however, quite debated until an unambiguous association was revealed between GRB 030329 and SN 2003dh (Hjorth et al., 2003; Stanek et al., 2003). SN 2003dh also showed properties almost identical to that of SN 1998bw. A third event has now filled in the picture with the spectroscopically confirmed association between GRB 031203 and SN 2003lw (Thomsen et al., 2004; Malesani et al., 2004). Because the detection of GRBs in $`\gamma `$-rays is unaffected by intervening gas and dust, they provide a powerful and possibly unbiased tracer of star-formation in the high-$`z`$ universe. This highlights the importance of studies of GRB host galaxies. Previous studies indicate that the majority of GRB host galaxies are blue and sub-luminous (e.g., Fruchter et al., 1999; Le Floc’h et al., 2003; Jakobsson et al., 2005). Based on these properties of GRB host galaxies, Watson et al. (2004) discussed whether a sizable portion of global star-formation occurs in small and rather unobscured, modestly star-forming galaxies that are too faint to appear in other surveys of star-formation activity, or whether GRBs trace only a fraction of the star-forming population, for example due to metallicity effects (e.g., Fynbo et al., 2003). The most nearby GRB host galaxies, and the ones that are undoubtedly directly connected to the deaths of massive stars, allow for a more detailed analysis. Watson et al. (2004) used X-ray observations to constrain the star-formation rates in the host galaxies of the three SN-GRB associations mentioned above. In the present study, we will instead use optical observations to address the same questions. In Sect. 2 we present the data on the hosts and shortly describe the performed data reductions. In Sect. 3 we characterize the host galaxies in terms of emission line fluxes, constraints on metallicities and extinctions, star formation rates, and physical dimensions. In Sect. 4 we discuss and compare the obtained results of the nearby GRB host galaxies with a population of blue compact galaxies. Our findings are summarized in Sect. 5 A cosmology where $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ and $`\mathrm{\Omega }_\mathrm{m}=0.3`$ is assumed throughout. The redshifts of the three host galaxies; SN 1998bw at $`z=0.0085`$, SN 2003lw at 0.1055 and SN 2003dh at 0.1685 then correspond to luminosity distances of 37, 487 and 810 Mpc, respectively. ## 2 Observations and data reduction ### 2.1 Spectroscopic data for the GRB 980425 host The host of GRB 980425 and SN 1998bw, ESO 184-G82, was extensively observed as part of several monitoring programmes for SN 1998bw. The evolution of the SN is described in detail elsewhere (Sollerman et al., 2000; Patat et al., 2001; Sollerman et al., 2002). Here we will make use of some of these observations to characterize the host galaxy. We have used the late spectroscopy from 13 June 1999 obtained at the ESO Very Large Telescope (VLT), using the Focal Reducer and low dispersion Spectrograph (FORS1) instrumentthanks: http://www.eso.org/instruments/fors1/ . These data are presented in Sollerman et al. (2000) and the spectra clearly reveal a number of narrow emission lines from the underlying host galaxy. The spectra were obtained using the 300V and 300I grisms and a $`1.^{\prime \prime }0`$ wide slit and cover the wavelength region $``$3700–9700 Å. The spectra were reduced in a standard way, including bias subtraction, flat fielding, and wavelength calibration using spectra of a Helium-Argon lamp. Flux calibration was done relative to the spectrophotometric standard star LTT 7379 (Hamuy et al., 1994). We did not only use the reduction of Sollerman et al. (2000), which aimed at reducing the presence of host contamination in the SN spectrum. Instead we re-reduced the data to better reveal the underlying emission line region. We also extracted spectra of two other H II regions in the part of the host covered by the slit (see Fig. 1). The absolute flux calibration of the combined spectrum was obtained relative to our well-calibrated broad-band SN photometry. However, since the slit does not cover the entire galaxy, we can not derive the total emission line flux and hence not estimate the global star formation rate. This is instead done from narrow band imaging (Sect. 2.2). ### 2.2 Narrow band imaging for the GRB 980425 host For ESO 184-G82 we have also retrieved narrow band imaging data from the ESO archive. In particular, we wanted to use the H$`\alpha `$ imaging to estimate the global star formation rate of the galaxy. These observations were performed on 3 August 2000 with the FORS1 instrument on the VLT. Three 5 minute exposures were obtained in a narrow filter centered on H$`\alpha `$ at the host galaxy redshift. Another five minute exposure was obtained in a zero-velocity narrow H$`\alpha `$ filter, to allow for continuum subtraction. The data were reduced in a standard way and the continuum flux-calibration was performed relative to a photometric standard star (PG 1657, Landolt, 1992) observed immediately before the galaxy at a similar airmass. The broad-band zero-points obtained from this standard star indicates that the night was photometric. ### 2.3 Broad band imaging for the GRB 980425 host We have re-analysed the late time $`BVRI`$ imaging of the host galaxy of SN 1998bw, which were originally used as templates for the SN template subtraction photometry in Sollerman et al. (2002). These are very deep images obtained with FORS1 at a time when the SN flux was negligible. We use it here to construct a spectral energy distribution (SED) for the galaxy. The photometry was calibrated against local photometric standard stars (Sollerman et al., 2002). The photometry for the entire galaxy is consistent with the photographic ESO-Uppsala catalogue (Lauberts & Valentijn, 1989). We also did photometry of the small H II region that contained the SN (see e.g., Fynbo et al., 2000; Sollerman et al., 2002). The magnitude was measured in an aperture with a diameter of 9 pixels (1$`.^{\prime \prime }`$8), and the local background was subtracted. ### 2.4 Spectroscopic data for the GRB 030329 host Spectroscopy of the host of GRB 030329 was obtained on 19 June 2003 with the FORS2 instrument on the VLT. As outlined by Gorosabel et al. (2005) these observations were conducted with the 300V grism and order sorting filter GG375, which efficiently covers the wavelength range from $``$3800–8800 Å. The 1$`.^{\prime \prime }`$3 wide slit was used and the seeing during the observations was below 0$`.^{\prime \prime }`$6. Given the small size of the object (see Fig. 1), the bulk of the flux from the host galaxy should be included in the slit. Also note that we used a position angle of 123.6 degrees East of North, which is well aligned with the orientation of the host galaxy (Fig. 1). The spectra were reduced in the standard way. Flux calibration was done relative to the spectrophotometric standard star Feige 67 (Oke, 1990). The absolute flux calibration of the combined spectrum was again obtained relative to broad-band photometry. A figure showing this spectrum is shown by Gorosabel et al. (2005). ### 2.5 High resolution spectra for the GRB 030329 host We retrieved early high-resolution spectra of GRB 030329 from the ESO archive, mainly with the aim of searching for the absorption lines of Na I D and hence constrain the amount of extinction along the line-of-sight to the burst. These lines are often used in SN studies for this purpose (see e.g., Sollerman et al., 2005). The observations were obtained 16 hours after the burst on 2003 March 30 with the Ultraviolet and Visual Echelle Spectrograph (UVES)thanks: www.eso.org/instruments/UVES/ on the VLT (Greiner et al., 2003). We reduced the spectra using the UVES-pipelinethanks: www.eso.org/observing/dfo/quality/ (vers. 2.0) as implemented in $`\mathrm{𝙼𝙸𝙳𝙰𝚂}`$. This reduction package allows for bias subtraction and flat-fielding of the data using calibration frames obtained during day time. Very accurate wavelength calibration was secured by comparison to ThAr arc lamp spectra. From the reduced spectrum we also confirm the redshift reported by Greiner et al. (2003) (see also Stanek et al., 2003; Hjorth et al., 2003), $`z=0.1685`$. The \[O III\] $`\lambda 5006.9`$ line was measured at $`\lambda 5850.69`$, which yields z=0.168525, and H$`\alpha `$ $`\lambda 6562.8`$ was measured at $`\lambda 7668.82`$ which gives z=0.168528. The measured wavelengths are corrected for the barycentric velocity. ## 3 Results ### 3.1 Emission line fluxes The measured emission line fluxes from the host galaxies of GRB 980425 and GRB 030329 are given in Table 1. In that table we also provide line fluxes for the host galaxy of GRB 031203, from the thorough analysis presented by Prochaska et al. (2004). We note that for GRB 030329 the main difference between the results presented by Gorosabel et al. (2005) and those by Hjorth et al. (2003) is the absence of significant \[N II\] emission. The detection of this emission line by Hjorth et al. (2003) in the spectrum from 1 May 2003 appears to have been spurious. It is not seen in any of the other spectra published by Hjorth et al. (2003) and the lack of \[N II\] emission is also noted by Matheson et al. (2003). ### 3.2 Metallicities From the measured emission line fluxes we have made an attempt to constrain the metallicities of the three host galaxies. We followed the methods outlined by Lee et al. (2003) and by Kewley & Dopita (2002) to derive metallicity estimates from the strong emission lines. These estimates are based on the R23 technique which is an empirical relation between the oxygen abundance and the intensity ratio of the strong oxygen emission lines (\[O II\] 3727 Å, and \[O III\] 4959,5007 Å) to H$`\beta `$. This relationship is, however, not unique and therefore provides two possible solutions (lower and upper branch) for the nebular oxygen abundance. This degeneracy may be resolved by including other lines such as \[N II\] 6584 Å. #### 3.2.1 The GRB 030329 host Using the prescription of Lee et al. (2003) for the emission line strengths of oxygen measured for the host of GRB 030329 we can derive a metallicity of either $`7.9`$ (lower branch) or $`8.6`$ (upper branch). The lower branch value would suggest a metallicity significantly below solar which would also be consistent with the luminosity of the galaxy (see, e.g., Fig. 6 in Lee et al., 2003) (or Fig. 4 in Lee et al., 2004). However, based on the emission lines alone we can not exclude a metallicity closer to solar (12+log(O/H)=8.7, Allende Prieto et al., 2001). Despite the depth of the VLT/FORS spectrum, \[O III\] 4363 Å and \[N II\] 6584 Å were not detected. Using the various diagnostic tools summarized by Kewley & Dopita (2002) we did not manage to break the degeneracy of the oxygen abundance solution based on the R23 ratio. In fact, using our upper limit on \[N II\] together with the \[O II\], H$`\alpha `$, and \[O III\] fluxes only gives weak constraints for most published diagnostic diagrams (Edmunds & Pagel, 1984; Kewley & Dopita, 2002), i.e., 12+log(O/H)$`<8.6`$. Also our \[S II\] detection gives only weak constraints when compared to \[N II\] and H$`\alpha `$. Hence we could not constrain the oxygen abundance of the host of GRB 030329 using the emission line fluxes, but it appears to be sub-solar. #### 3.2.2 The GRB 980425 host For the host galaxy of GRB 980425 we have performed a similar analysis based on the emission line fluxes. We caution that the absolute fluxing in the bluest part of the spectrum, where \[O II\] is located, may be affected by systematic uncertainties in the flux-calibration. This, as well as uncertainties in the extinction corrections may affect the metallicity estimates, but does not alter the overall results. We extracted spectra from three spatial locations within the host galaxy (Fig. 1). The metallicities were then derived using the recipe in Kewley & Dopita (2002). In this case, we have enough information to support the upper branch of the R23 diagnostic, and find the metallicities shown in Table 2. Thus, this galaxy does not have a very low metallicity. There may be a slight variation in the metallicity across the galaxy, but given the uncertainties we do not regard this as a significant result. We note that the emission lines of \[Ne III\] are stronger at the site of the burst. As argued by Bloom et al. (1998) this is indicative of a high degree of ionization, suggesting a substantial population of young and massive stars. Such an interpretation is also supported by our modeling of the broad-band spectral energy distribution (Sect. 3.5). ### 3.3 Extinction The collected dataset allows some estimates of the extinction in the hosts. The extinctions are also needed for estimating the star formation rates below. For the GRB 030329 host galaxy there are several estimates of the extinction available. First, the photometric spectral energy distribution of the host as fitted by Gorosabel et al. (2005) favours a starburst galaxy with an intrinsic $`E(BV)0.2`$ mag. The Galactic extinction in this direction is estimated to be $`E(BV)=0.025`$ mag (Schlegel et al., 1998). Matheson et al. (2003) also argued for a low extinction. They based their arguments on the Balmer decrement and also made an estimate based on the assumed power-law properties of the afterglow emission. The latter method gave a limit on the extinction towards the burst of $`E(BV)=0.04\pm 0.08`$, implying that there is no evidence for extragalactic dust along the line of sight between us and GRB 030329. The high-resolution spectroscopy of the afterglow of GRB 030329 discussed above (Sect. 2.5) also speaks in favor of a very low amount of dust towards the GRB. At the position of the Galactic sodium lines we do detect a weak line. The $`\lambda 5890`$ component has an equivalent width (EW) of $`40`$ mÅ. This small EW is fully consistent with the low amount of Galactic extinction in this direction (e.g., Hobbs, 1974; Sollerman et al., 2005). At the position of the Na I D lines at the redshift of GRB 030329 we detect no significant absorption. Any such absorption must be at least four times weaker than the Galactic component. Although the degree of ionization could be different in the GRB 030329 host and the Galaxy, this at least points to a very small amount of reddening along the line of sight to GRB 030329. We must remember, however, that the line-of-sight towards the GRB need not be representative for the entire host galaxy. It may even be that dust is destroyed by the burst itself (e.g., Waxman & Draine, 2000; Galama & Wijers, 2001; Fruchter et al., 2001). For the GRB 980425 host galaxy we have limited information on the extinction. The Galactic extinction in this direction is estimated to be $`E(BV)=0.059`$ mag (Schlegel et al., 1998). According to Patat et al. (2001) there was no sign of Na I D absorption in high-resolution spectra obtained at the SN maximum. These authors used this finding to argue for an extinction of $`E(BV)<0.065`$ mag towards the SN, i.e., a again very small amount of reddening. Also the Balmer decrement from the spectra indicate a low reddening, which is also consistent with our SED modeling (Sect. 3.5.1). ### 3.4 Star Formation Rates There are several ways to estimate star formation rates (SFR) in galaxies. One option is simply to use the emission line luminosities of H$`\alpha `$ or \[O II\] and follow the prescription by Kennicutt (1998). For the GRB 980425 host galaxy we have used narrow band imaging to estimate the flux in H$`\alpha `$, since this galaxy is much larger than the extent of our spectroscopic slit. The final continuum subtracted image is shown in Fig. 1. The host galaxy is clearly visible. We also note that no other source was detected in the field of view. From this image we derived an integrated H$`\alpha `$ flux of $`2.6\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. The integrated EW is 60.6 Å. If we correct these for a 10$`\%`$ contribution from the \[N II\] emission line within the filter band (as estimated from the spectroscopy), and for a Galactic extinction of $`E(BV)=0.059`$ mag we get for the distance of GRB 980425 a luminosity of L$`{}_{\mathrm{H}\alpha }{}^{}=4.4\times 10^{40}`$ erg s<sup>-1</sup>. The corrected EW is 54.5 Å. We thus derive a global SFR of $`0.35`$M yr<sup>-1</sup> while the bright H II region northwest of the explosion site of SN 1998BW has $`0.11`$M yr<sup>-1</sup>. These estimates have also been corrected for Galactic extinction and for a 10$`\%`$ contamination of \[N II\] emission inside the narrow band filter. For the GRB 030329 host we obtain the following integrated SFR; SFR$`{}_{\mathrm{H}\alpha }{}^{}=0.22`$M yr<sup>-1</sup> SFR$`{}_{[\mathrm{O}\mathrm{II}]}{}^{}=0.22`$M yr<sup>-1</sup> We note that these values are consistent with the values reported by Hjorth et al. (2003). The H$`\alpha `$ line sits close to a telluric absorption line and was not used by Gorosabel et al. (2005) to estimate the SFR. These values are for a Galactic extinction only. An extinction as suggested by Gorosabel et al. (2005), $`E(BV)=0.2`$ would alter these numbers to SFR$`{}_{\mathrm{H}\alpha }{}^{}=0.32`$M yr<sup>-1</sup> SFR$`{}_{[\mathrm{O}\mathrm{II}]}{}^{}=0.48`$M yr<sup>-1</sup> This of course applies only to the part of the galaxy included in the slit - although we argue that this was actually the major part. These estimates are similar to the SFR reported by Matheson et al. (2003), 0.5 M yr<sup>-1</sup>, although the agreement appears to be somewhat accidental. We were unable to reproduce their high L also from their publicly available spectra, but on the other hand they assume a much lower extinction correction. The SFR estimated from the optical emission lines must be regarded as lower limits, since there may also exists an extinguished star-forming population, as in many other starburts. For example, the host of GRB 000210 revealed a SFR $`23`$M yr<sup>-1</sup> from the optical emission lines (Gorosabel et al., 2003), while a tentative sub-mm detection implies a star formation of several hundred solar-masses per year (Berger et al., 2003). In the case of the SN/GRB hosts we also have upper limits on the star formation as obtained from the unbiased X-ray view exploited by Watson et al. (2004). For the GRB 980425 host they find that the X-ray flux within the optical extent of this galaxy is entirely dominated by two point sources $`1.5`$<sup>′′</sup> apart, one of which is coincident with the radio position of SN 1998bw and is almost certainly associated with it (Watson et al., 2004; Kouveliotou et al., 2004). From the X-ray emission they estimated a total SFR of $`2.8\pm 1.9`$ M yr<sup>-1</sup>. Together with our estimate of 0.35 M yr<sup>-1</sup> the range is therefore quite well constrained for this host galaxy. For GRB 030329 we have estimated SFR$`0.4`$ M yr<sup>-1</sup> while Watson et al. (2004) estimate a SFR of massive stars (M$`>5M_{}`$, see Watson et al. 2004; Grimm et al. 2003) of less than $`31\pm 13`$ M yr<sup>-1</sup>. Using a salpeter IMF (Salpeter, 1955; Persic et al., 2004; Watson et al., 2004) this corresponds to a total SFR of $`<200\pm 80`$ M yr<sup>-1</sup>. ### 3.5 Luminosities, Physical sizes and Galaxy types. #### 3.5.1 The GRB 980425 host The host of GRB 980425, ESO 184-G82, is nearby enough to be well resolved with ground based telescopes (Fig. 1) and appears to be a late type spiral with a bar (SBc). The beautiful HST image displayed by Fynbo et al. (2000) shows the optical appearance of the galaxy to be dominated by a large number of high surface brightness starforming regions, especially in the southern spiral arm where the GRB/SN occurred. The total BVRI magnitudes we obtained from our VLT images down to an isophote of 25 magnitudes arcsec<sup>-2</sup> are reported in Table 3. For the adopted distance this gives an absolute magnitude of M$`{}_{B}{}^{}=17.65`$. Adopting M$`{}_{B}{}^{}=21`$ we thus find that this galaxy has $`L=0.05L^{}`$. We can use this and the above derived SFR to also get the specific star formation rate (SSFR). For GRB 980425 we thus derive a SSFR$`7`$M yr<sup>-1</sup> (L/L)<sup>-1</sup>. The major axis diameter of the galaxy is 67 arcseconds at the B=26.5 mag arcsec<sup>-2</sup> isophote, corrected for foreground Milky way extinction, i.e. this is the Holmberg diameter. The minor axis is 57 arcseconds. At a distance of 37 Mpc this corresponds to a physical size of $`12\times 10`$ kpc. Comparing our broadband magnitudes to empirical galaxy colours (e.g., Coleman et al., 1980) we see that the host appears to be a typical, subluminous late type spiral. To quantitatively compare the colors of the host galaxy we calculated a set of models using the code PEGASE.2 (Fioc & Rocca-Volmerange, 1997, 1999) which includes both stellar and nebular emission. This was done using the actual filter profiles and CCD sensitivity for FORS1, and all models assume a Salpeter IMF in the mass interval 0.1-120 M. Using a range of different metallicities and star formation timescales we could then compare these models to our data using least-square fitting. The entire galaxy is well fit with a continuous star formation history. Both an e-folding time of 3 Gyr and 15 Gyr give good agreements without additional extinction. The results are not very sensitive to metallicity. Independently, the H$`\alpha `$ EW is well fit by such a scenario. Both these scenarios are also able to reproduce the integrated luminosity of the galaxy given that the current SFR has operated for 5-7 Gyrs. The galaxy can thus not be regarded as a starburst galaxy. #### 3.5.2 The GRB 980425 progenitor mass For the local SN/GRB environment we have conducted a similar exercise. The photometry is presented in Table 3 and the results of the modeling are quantified and summarized in Table 4. We assumed an instantaneous burst and the same IMF parameters as above. The best fitting age is for each metallicity a well defined minimum, where a change in age as small as $`\pm 1`$ Myr typically increases the RMS deviation with a factor two or more. For all models, the best-fitting internal reddening is found to be $`E(BV)=0.05`$. As is illustrated in Table 4, a change in $`E(BV)`$ of $`\pm 0.05`$, leads to an increased RMS with typically 50 to $`100\%`$, but does not affect the best fitting age. Hence, the uncertainties on the derived ages are very small ($``$1 Myr) in a statistical sense. If we take the best fitting age as the lifetime of the supernova progenitor we can estimate its mass from comparison with stellar tracks (Bressan et al., 1993; Fagotto et al., 1994; Meynet et al., 1994). This gives a ZAMS mass of 30$`\pm `$5 M, in accordance with most models for Type Ic SNe in general, and for collapsars in particular. #### 3.5.3 The GRB 030329 host For GRB 030329 the magnitudes reported by Gorosabel et al. (2005) imply $`L=0.016L^{}`$. With our measured (extinction corrected) SFR this gives a SSFR of $`25`$M yr<sup>-1</sup> (L/L)<sup>-1</sup>. As noted by Gorosabel et al. (2005) this is a very high SSFR compared to most galaxies in the Hubble Deep Field. To constrain the physical size of this host galaxy we used images obtained with the ACS onboard the Hubble Space Telescopethanks: programme 9405; P.I. A. Fruchter. . We measured the extent of the host on ACS images obtained on 2004 May 24, i.e., 422 days after the GRB and thus long after the afterglow contribution had vanished. The FWHM of the host is about $`530\times 930`$ pc in the F435W filter image. The F606W and F814W filters show a slightly larger FWHM of $`620\times 1030`$ pc. This is similar to the radius estimated by fitting a Sersic model to the surface brightness profile (e.g., Warren et al., 2001), which gave a radius of 0.26 arcseconds corresponding to 750 pc for the F606W filter. We also estimated the Holmberg diameter to compare directly with the estimate of the GRB 980425 host galaxy. This gave instead $`1.^{\prime \prime }4`$, corresponding to 3.9 kpc. The main conclusion of this exercise is that the host galaxy of GRB 030329 is a very compact galaxy, with both an absolute magnitude and extension similar to that of the SMC. Furthermore, our UVES spectroscopy (Sect. 2.5) resolved the H$`\alpha `$ line with a FWHM of $`55`$ km s<sup>-1</sup>, which for a radius of 0.75 kpc corresponds to a dynamical mass of $`5\times 10^8`$ M (e.g., Östlin et al., 2001). This may be an underestimate if the H$`\alpha `$ flux is dominated by a central burst and do not trace the full potential well. The dynamical mass is similar to the mass of SMC. In Fig. 1 we show the host galaxies of GRBs 030329 and 980425 on the same physical scale. This comparison reveals that these host galaxies are not that different. The VLT images of GRB 980425 are considerably deeper and surface brightness dimming will also suppress the fainter structures in the more remote galaxy. One difference is the location of the GRB within the galaxy. GRB 030329 appears to have exploded right on the brightest pixel in the host galaxy. The location is marked by a ring and a cross in Fig. 1. The occurrence of GRB 980425 is instead way outside the center of the galaxy, at a projected distance of about 2.2 kpc. This region does not appear to be special in any way, compared for example to the very actively star forming region 850 pc northwest of the explosion site (Fig. 1). We note that at cosmological distances such a spatial difference would not have been possible to resolve. From Gorosabel et al. (2005) we know that the host of GRB 030329 has a SED that is best fit by a starburst galaxy template. We have also measured the photometry on the HST images and found m(F435W)=23.29, m(F606W)=23.00 and m(F814W)=22.83. These are aperture corrected AB magnitudes. The errors are estimated to be about 0.03 magnitudes. The B and V band results are fairly consistent with Gorosabel et al. (2005) although the I band is measured to be fainter in the HST images. Looking also at the spatially resolved photometry, we note that the colors are significantly redder in the outer parts of the galaxy. This means that using an integrated magnitude will overestimate the age of the stellar population (e.g., Östlin et al., 2001). This would make the ages estimated by Gorosabel et al. (2005) more consistent with the expected ages from very massive GRB progenitors. ## 4 Comparisons Having established and collected the properties of the host galaxies of GRB 980425 and GRB 030329 we will in this section compare these properties with those of other host galaxies. First, we compile the properties of the third nearby GRB/SN host galaxy, that of GRB 031203, in the same way as we have done for the other two host galaxies. Thereafter we compare this small sample of galaxies with a sample of more distant Blue Compact Galaxies. ### 4.1 Comparison to the GRB 031203 host The third nearby SN/GRB host galaxy, the host of GRB 031203, has been spectroscopically studied in detail by Prochaska et al. (2004). Using their published reddening corrected emission line fluxes we performed the same kind of emission line analysis as for the other hosts. The resulting metallicities are given in Table 2. We note that our values are somewhat higher than those found by Prochaska et al. (2004, 12+log(O/H)=8.0). This is probably because the method we have used does not take into account the additional temperature information from the fainter \[O III\] 4363 Å line. This makes our abundance determinations less secure, and is included here to indicate the uncertainties in the metallicities also for the other two host galaxies. When it comes to extinction, Prochaska et al. (2004) estimated this from the Balmer decrement to be $`E(BV)=1.17`$. Most of this, $``$1.04 mag, is due to extinction in our own galaxy (Schlegel et al., 1998). Hence the internal reddening also in this host galaxy is likely to be modest. Given this extinction, we take the line fluxes from Prochaska et al. (2004) and use the same method as above to infer SFR$`{}_{\mathrm{H}\alpha }{}^{}=13.2`$M yr<sup>-1</sup>, SFR$`{}_{[\mathrm{O}\mathrm{II}]}{}^{}=9.6`$M yr<sup>-1</sup>. The X-ray data (Watson et al., 2004) suggest a massive SFR of at most $`24\pm 17`$ M yr<sup>-1</sup>, corresponding to a total SFR of $`<150\pm 110`$ M yr<sup>-1</sup>. The low K-band luminosity (L$``$L$`{}_{K}{}^{}{}_{}{}^{}`$/5) and the extrapolated B-band magnitude (Prochaska et al., 2004) corresponds to $`L=0.28L^{}`$. This gives a SSFR=$`39`$ M yr<sup>-1</sup> (L/L)<sup>-1</sup>. For this host galaxy we have limited information about the size and galaxy type. Using an $`I`$-band image obtained at the Danish 1.54 m telescope 46 days past explosion (Thomsen et al., 2004) we compared Gaussian fits to the host galaxy with fits of the nearby stars. The galaxy is clearly not a point source, but with a seeing of $`1.^{\prime \prime }0`$ it is only marginally resolved. Taking into account the extent of the point spread function this means that the galaxy FWHM is about 1 – 2 kpc on the sky. This is slightly larger than the very compact GRB 030329 host, but not by much (Fig. 1). A J-band image presented by Gal-Yam et al. (2004) shows that the GRB occurred in the central regions of this galaxy. ### 4.2 Comparison to other galaxies The properties of the three SN-GRB host galaxies are summarized in Table 5. The star formation rates given in the table are averages between those obtained for H$`\alpha `$ and \[O II\], when available. With only three firmly established GRB/SN host galaxies any conclusion regarding this population will be tentative, but it seems that their overall properties are consistent with those established for more distant GRB host galaxies; small, sub-luminous, metal-poor blue galaxies (e.g., Fruchter et al., 1999; Le Floc’h et al., 2003; Fynbo et al., 2003; Christensen et al., 2004). The host galaxy of GRB 980425 has, however, not a very low metal abundance - and should not be classified as a starburst galaxy. Christensen et al. (2004) compared the specific star formation rates (UV to optical colours, since the SFRs were estimated from the near-UV continuum) of GRB hosts to galaxies in the Hubble deep field, and concluded that the GRB hosts have higher SSFR than $`95\%`$ of the HDF galaxies. Such a comparison may give the appearance that the GRB hosts have very special properties. However, the HDF comprises a plethora of galaxies of different types with different distances and selection functions. Given that the hosts under study in the present paper are rather nearby, it would be interesting to compare these with galaxies in the local universe with similar properties. Although not comprising a very homogeneous class, Blue Compact (dwarf) Galaxies (BCGs) and low luminosity emission line selected galaxies (H II-galaxies) typically have luminosities below L, low metallicity (Izotov & Thuan, 1999), low extinction (Izotov et al., 1997) and active star formation (Kunth & Östlin, 2000). These are properties often ascribed to GRB hosts and the BGC population may therefore serve as an interesting local comparison sample. Gil de Paz et al. (2003) present visual and H$`\alpha `$ luminosities for a large sample of local BCGs, which we have converted into specific star formation rates. We also added a few luminous BCGs with very active star formation from Östlin et al. (2001) to this sample, with the SFRs recalibrated by the Kennicutt (1998) SFR–H$`\alpha `$ relation for the sake of homogeneity. In Fig. 2 we compare the star formation rates and luminosities of the combined BCG sample with that of the three SN-GRB hosts of this paper. For this comparison we used the SFRs inferred from H$`\alpha `$ only, and since Gil de Paz et al. (2003) do not present internal extinction, we use vales corrected for Galactic extinction only in order to be consistent. The combined samples cover a range in SSFR from $`1`$ to 100 M yr<sup>-1</sup> (L/L)<sup>-1</sup> while the typical BCG has a SSFR close to 10. While this large range of values confirms that this class of galaxies is not very homogeneous, we see that the SN/GRB host galaxies have similar properties to this class. It is among these galaxies that we find the most active star forming sub-luminous galaxies in the local universe. The BCGs with SSFR $`<10`$ have moderately active star formation and this would also be the location occupied by ordinary dwarf irregulars and disk galaxies, such as the host of GRB 980425. The other two SN/GRB hosts would probably have been classified as actively starforming BCGs if they would have been more nearby. BCGs are believed to account for a significant amount of the SFR also at higher redshifts. Still, in a magnitude limited sample, such galaxies are in minority. In this respect the results by Gorosabel et al. (2005) and Christensen et al. (2004) can be understood – the majority of the (SN/)GRBs occur in the few percent of the magnitude limited population that has the highest SSFR, both at low and high redshifts, i.e., galaxies with properties similar to BCGs/H II-galaxies (see also Courty et al., 2004). ## 5 Conclusions We have mined several archives and data sets in order to characterize the host galaxies that are known to have harbored a GRB-SN. The picture of compact low luminosity, metal poorish galaxies in a starforming phase is established, consistent with other studies of ordinary and more distant GRBs (Le Floc’h et al., 2003; Christensen et al., 2004). The local extinction in these galaxies appears to be small. We have proposed that a population of Blue Compact Galaxies have similar properties to the GRB-SN host galaxies. For the host of GRB 980425 we have for the first time derived the star-formation rate. We have shown that this galaxy is not very metal-poor, and that a population study based on the broad-band photometry is consistent with a normal starforming galaxy with continuous star formation over 5-7 Gyrs, i.e., not with a starburst galaxy. A similar investigation of the spatially resolved H II region where the GRB occurred gives us an estimate of the GRB progenitor mass of $`>30M_{}`$. This is consistent with theoretical scenarios of SN 1998bw being due to a very massive star (e.g., Iwamoto et al., 1998). Today, collapsar models - involving the collapse of massive stars - are the favoured models for long GRBs (e.g., MacFadyen & Woosley, 1999; Hjorth et al., 2003), and it is therefore very encouraging that this population study confirms such models. Similar observational contraints on the mass of supernova progenitors are becoming very useful for constraining SN models (see e.g., Maund & Smartt, 2005, and references therein). For the host of GRB 030329 we provide high-resolution HST imaging which reveals the morphology and location of the burst. It shows the host to be a very compact galaxy indeed. An early high-resolution spectrum provides a very accurate redshift determination and the lack of sodium lines in this spectra supports a very low extinction towards the GRB. The width of the resolved H$`\alpha `$ provides an estimate of the dynamical mass of the galaxy. Our low-resolution spectra give useful constraints on the metallicity, but we also show that the analysis is less straight-forward than previously acknowledged. This study thus provides a smorgasbord of what can be learned from optical data of relatively nearby host galaxies. We hope that more such galaxies will be discovered with the Swift satellite. To characterize this population is important to understand selection biases in determining the high-z starformation rates via GRB selected galaxies. Swift has the potential to point us to a large number of such distant galaxies. The datasets used in this investigation were mainly obtained to study the individual SNe-GRBs and were not really optimized for the host studies. This is particularly true for the spectroscopy. Clearly, better observations are possible to obtain when the SNe have faded away. Acknowledgements. We want to thank Javier Gorosabel for discussions about several important aspects of this paper. We also want to thank Lisa Kewley for providing her metallicity script. Thanks to Christina Thöne for comments. Part of this research was conducted at the Dark Cosmology Centre funded by The Danish National Research Foundation. Important VLT observations were conducted on ESO Director’s Discretionary Time and were obtained by Paranal staff. We are grateful for these efforts.
warning/0506/cond-mat0506143.html	ar5iv	text	# Evidence of correlation in spin excitations of few-electron quantum dots ## Abstract We report inelastic light scattering measurements of spin and charge excitations in nanofabricated AlGaAs/GaAs quantum dots with few electrons. A narrow spin excitation peak is observed and assigned to the intershell triplet-to-singlet monopole mode of dots with four electrons. Configuration-interaction theory provides precise quantitative interpretations that uncover large correlation effects that are comparable to exchange Coulomb interactions. Electrons confined to semiconductor quantum dots (QDs) have novel ground and excited states that manifest Coulomb interactions at the nanoscale Reimann . States of very few electrons are prime candidates for spintronic applications and for the implementation of quantum bits in nanoscale devices loss98 . Great attention is therefore devoted to the study of spin physics in the regime of few-electron occupation and to experimental methods capable of reading the state of spin in the QD kouw04 . The interpretation of experiments on few-electron QDs often requires descriptions beyond mean-field, such as Hartree-Fock (HF) Reimann ; haw04 . In addition to their relevance for quantum information encoding, these correlated states have significant interest for the investigation of fundamental effects egger ; oliver . Transport in few-electron QDs coupled to leads and excitonic optical recombination measurements have explored exchange energies and spin relaxation times. These remarkable experiments offer evidence for roles played by interactions that emerge as the number of electrons in the QD is changed Taruch ; exchange ; kouwe01 ; fuji02 ; kogan03 ; kouw04 ; petta04 ; zhumbul04 ; haw04 ; bayer98 ; fin04 . In this Letter we report resonant inelastic light scattering experiments in low electron density GaAs/AlGaAs QDs that probe low-lying neutral excitations. These are inter-shell monopole excitations (with change in Fock-Darwin (FD) shell and without change in angular momentum, as required by light scattering selection rules lockwood ; schuller ). We detected two broad inter-shell modes that we interpret as excitations of electrons from the two populated lowest shells. Each of these two modes is split by exchange and depolarization effects into a $`\mathrm{\Delta }S=1`$ (spin) and a $`\mathrm{\Delta }S=0`$ (charge) excitation lip ; delg , where $`\mathrm{\Delta }S`$ represents the change of the total spin of the QD associated to the electronic mode. A prominent feature of the spectra is an additional mode peculiar to the regime of few-electron occupation that emerges at low temperature and low excitation power. It occurs as a very narrow peak with light scattering polarization selection rules for spin excitations and is interpreted as a $`\mathrm{\Delta }S=1`$ intershell spin mode characteristic of a $`S=1`$ triplet ground state with four electrons. We argue that the observed splitting between the $`\mathrm{\Delta }S=1`$ and $`\mathrm{\Delta }S=1`$ spin modes represents a direct manifestation of the role of interactions in the excitation spectra of few electron QDs. Numerical evaluations within a configuration-interaction (CI) theory ront04 support the interpretation that links the new spin mode to the triplet-to-singlet (TS) inter-shell excitation of a QD with four electrons and offer quantitative insights on the role of interactions in this regime. The CI evaluations reproduce the experimental light scattering spectra with a great precision that is not achieved by HF theory. Comparisons of mean field and CI calculations uncover large exchange and correlation terms of electron interactions that in the case of the four-electron triplet state are found to be comparable to quantum confinement energies. Samples were fabricated from a 25 nm wide, one-side modulation-doped Al<sub>0.1</sub>Ga<sub>0.9</sub>As/GaAs quantum well with measured low-temperature electron density $`n_e=1.1\times 10^{11}`$ cm<sup>-2</sup> and mobility of $`2.7\times 10^6`$ cm<sup>2</sup>/Vs. QDs were produced by inductive coupled plasma reactive ion etching. QD arrays (with sizes $`100\times 100`$ $`\mu `$m containing $`10^4`$ single QD replica) were defined by electron beam lithography with different diameters. Deep etching (below the doping layer) was then achieved. Here we focus on QDs having lateral lithographically-defined diameters of 210 nm (shown in Fig. 1 side panels) that we expect to be close to the regime of full electron depletion kir02 . The experiments were performed in a backscattering configuration ($`q2\times 10^4`$ cm<sup>-1</sup> where $`q`$ is the wave-vector transferred into the lateral dimension) with temperatures down to $`T=1.9`$ K. A tunable ring-etalon Ti:sapphire laser was focussed on 100 $`\mu `$m-diameter area and the scattered light was collected into a triple grating spectrometer with CCD multichannel detection. A convenient description of single-particle QD levels is provided by FD orbitals Reimann with energies given by $`\epsilon _{nm}=\mathrm{}\omega _0(2n+\|m\|+1)`$, where $`n=0,1,\mathrm{}`$, $`m=0,\pm 1,\mathrm{}`$ are the radial and azimuthal quantum numbers, respectively, and $`\mathrm{}\omega _0`$ is the harmonic confinement energy. The FD shells are defined by an integer value of $`N_{\text{shell}}=2n+\|m\|`$ with well defined atomic-like parity. QD states can be classified in terms of the $`z`$-component total angular momentum $`M`$, total spin $`S`$, and its $`z`$-component $`S_z`$. Selection rules in QDs dictate that the monopole transitions with $`\mathrm{\Delta }M=0`$ ($`\mathrm{\Delta }N_{\text{shell}}=2,4,\mathrm{}`$) are the inter-shell modes active in light scattering experiments strenz ; schuller ; lockwood . This non-interacting picture of intershell transitions is shown in the left part of Fig. 4, where the lowest energy dipole ($`\mathrm{\Delta }N_{\text{shell}}=1`$) and monopole ($`\mathrm{\Delta }N_{\text{shell}}=2`$) modes are represented. Figure 1(a) shows representative low-temperature spectra of inter-shell spin and charge excitations that are detected with crossed and parallel polarization between incident and scattered light, respectively vit01 . Pairs of peaks are seen at energies close to $`4`$ meV and $`78`$ meV and interpreted as monopole excitations with $`\mathrm{\Delta }S=1`$ (spin) or $`\mathrm{\Delta }S=0`$ (charge). In the non-interacting FD picture these two excitations are degenerate but they split in the presence of exchange and depolarization contributions. A characteristic feature of these doublets is their significant linewidth that is attributed largely to inhomogeneous broadening due to the distribution of electron occupations of the dots as described below. The spectra in Fig. 1(a) reveal a much sharper excitation in the spin channel at $`5.5`$ meV. To interpret the origin of this sharp spin mode we note that an additional TS intershell spin mode with $`\mathrm{\Delta }S=1`$ can occur if the ground state is a triplet with $`S=1`$ and the excited state is a $`S=0`$ singlet state (see the calculated levels shown in Fig. 4). Such TS excitation is split from the $`\mathrm{\Delta }S=+1`$ mode seen at lower energy by the difference in exchange and correlation contributions. On this basis we identify the sharp peak at $`5.5`$ meV with the TS ($`\mathrm{\Delta }S=1`$) intershell spin excitation. According to Hund’s rules a triplet ground state occurs only when two electrons are in a partially populated shell as is the case of QDs with four electrons. This is confirmed by the calculations described below. In this interpretation the narrow width is simply explained as due to the absence of inhomogeneous broadening from the distribution of the electron population of the QDs. We used the full CI approach ront04 ; Haw95 ; method for the numerical evaluation of the energy and intensity of low-lying spin and charge excitations of the interacting system with $`N`$ electrons. The correlated wavefunctions of ground and excited states are written as superpositions of SDs, $`\|\mathrm{\Phi }_{\{\alpha _i\}}=_{i=1}^Nc_{\alpha _i}^{}\|0`$, obtained by filling in the single-particle spin-orbitals $`\alpha `$ with the $`N`$ electrons in all possible ways \[$`c_\alpha ^{}`$ creates an electron in level $`\alpha (n,m,\text{or})`$\]. The resulting Hamiltonian is first block diagonalized, fully exploiting symmetries method . Finally, the Hamiltonian is diagonalized via Lanczos method in each $`(M,S,S_z)`$ sector, giving the low-energy excited states. The resonant Raman transition matrix elements $`M_{\text{FI}}`$ between the fully interacting ground and excited states $`\|\text{I}`$ and $`\|\text{F}`$, respectively, are obtained, after the CI calculation, from $`M_{\text{FI}}=_{\alpha \beta }\gamma _{\alpha \beta }\text{F}\|c_\alpha ^{}c_\beta \|\text{I}`$, where $`\gamma _{\alpha \beta }`$ is the two-photon process matrix element between $`\alpha `$ and $`\beta `$ spin-orbitals, as defined in Ref. Steinebach, within second order perturbation theory in the radiation field and containing resonant denominators. $`\gamma _{\alpha \beta }`$ causes the enhancement of the spectrum intensity when the laser energy resonates with the optical gap gamma . Figure 1(b) displays the calculated spectra for $`N=4`$ and $`\mathrm{}\omega _0=4`$ meV. The latter value was determined by fitting the peak energy position in the experimental spectra shown in Fig. 1(a). An independent check for this value of $`\omega _0`$ and $`N`$ comes from the empirical relation given in Ref. Reimann, \[Eq. (11)\] linking $`N`$, $`\omega _0`$, and the electron density, $`n_e`$, which gives $`n_e=1.2\times 10^{11}`$ cm<sup>-2</sup>, in good agreement with the experimental value. It can be seen that among all calculated excitations with $`\mathrm{\Delta }M=0`$, only a few of them turn out to have significant intensities, generating discrete spectrum lines (with a phenomenological broadening chosen to reproduce the measured TS linewidth) in very good agreement with the experimental ones. It can also be noted that more than one excitation gives a significant contribution to the spectra at energies above the TS mode. This is consistent with the observed larger linewidths for the higher-energy excitation pairs. Figure 2(a) reports the evolution of the calculated spectra as a function of $`N`$. As expected, the TS ($`\mathrm{\Delta }S=1`$) mode is peculiar to $`N=4`$ and it is not observed at any other explored electron occupation configurations. The excitations of Fig. 2 show a redshift of the lowest-energy features in both channels as $`N`$ (and $`n_e`$) increases brocke due to screening effect. Because of the exchange energy gain of excited states, the spin channel energy is systematically lower than the charge excitation energy. This large sensitivity of the light scattering spectra on particle occupation is at the origin of the difference between the observed linewidths of our inter-shell excitations. Comparing the evolution of peak energies shown in Fig. 2 with measured linewidths we conclude that a distribution of electron occupation between $`4`$ and $`6`$ characterizes our QD arrays. It also indicates that the light scattering method allows to probe excitation spectra of few-electron QDs with single-electron accuracy despite the relatively large number of QDs illuminated. Consistent with the assignment that links the $`\mathrm{\Delta }S=1`$ mode to those selected QDs that have four electrons, is the observed sharp linewidth of 0.4 meV which is much lower than the linewidths of the other spin and charge transitions. The evolution of the spin transitions at different incident laser intensities shown in Fig. 3(a) confirms that the QDs are in the few-electron regime. As the intensity increases we expect additional electrons to be photo-generated. Consistent with Fig. 2, we found that the peaks display a red shift and that the TS transition disappears at around I = 1 W/cm<sup>2</sup>, suggesting that at this intensity all the QDs have more than 4 electrons and therefore the number of those photo-generated is at least one. In addition, contrary to the other inter-shell modes the intensity of the TS spin excitation decreases significantly as temperature increases \[Fig. 3(b)\] with an estimated activation gap of $`0.7\pm 0.3`$ meV. At such low energy a possible thermally populated excited level is the singlet state without any change in orbital occupation. This energy thus provides an estimate of the low lying intra-shell singlet-triplet transition of the four-electron QDs and it compares well with CI estimate of 0.8 meV. A specific feature of the low-$`N`$ regime studied here is that states are represented by superpositions of many different SDs to incorporate both radial and angular spatial correlation Steinbach . The side diagrams of Fig. 4 represent the weighted SDs in the CI expansion of the $`N=4`$ ground and excited states involved in the three transitions indicated by arrows in Fig. 1(b). We depicted the states corresponding to the maximum allowed $`S_z`$, while in the actual calculation we only considered the degenerate states with $`S_z=0`$ method . Figure 4 also shows ground- and excited-state energies calculated with different approximations that provide evidence for correlation effects in the excitation spectra. In the FD picture the energy difference between consecutive levels is given by $`\mathrm{}\omega _0`$ = 4 meV. In the HF approach, spin-orbitals are computed self-consistently Rontani99 . The energy difference between the three spin configurations is given by bookkeeping the exchange energy $`K_{ab}`$ gained each time two electron spins, occupying any orbitals $`a`$ and $`b`$, are parallel to each other. This gain is accounted for by the Coulomb exchange integral between orbitals $`a`$ and $`b`$ described by FD wave functions. This approach neglects spatial correlation among electrons. Correlation effects are included in the CI approach, leading to the theoretical spectra in Figs. 1 and 2 and to the quantitative agreement with experiments shown in Fig. 4. The comparison between HF and CI (Fig. 4) suggests that correlation affects the relative splittings between excited states, even reversing their relative amplitudes: The $`S=1`$ state is nearer to $`S=0`$ than to $`S=2`$ in HF, while the opposite is true in CI, in agreement with the experiment. As suggested by the decreasing contribution of the most weighted SD configurations indicated on the right in Fig. 4, correlation effects are small for the ground and the $`S=2`$ excited state, but become increasingly important for excited states with smaller $`S`$. As $`S`$ decreases, exchange interaction is less effective in keeping electrons far apart and excited states become more correlated. Note that the relative amplitudes of the calculated HF and CI gaps are quite insensitive to the specific value of $`\mathrm{}\omega _0`$ and we found that the measured splittings among the spin modes can only be reproduced by CI calculations, no matter the value of $`\mathrm{}\omega _0`$. In conclusion, we reported inelastic light scattering measurements of spin transitions in nanofabricated quantum dots. The characteristic excitations of the triplet configuration with four electrons have been identified and theoretically evaluated. We have shown that light scattering methods offer a wealth of information on the physics of spin states in QDs with few electrons. We acknowledge support from the Italian Ministry of Foreign Affairs, Italian Ministry of Research (FIRB-RBAU01ZEML and COFIN-2003020984), CINECA-INFM Supercomputing Project 2005, European Community’s Human Potential Programme (HPRN-CT-2002-00291), National Science Foundation (DMR-03-52738), Department of Energy (DE-AIO2-04ER46133) and a research grant of the W. M. Keck Foundation. We are grateful to SENTECH-Berlin for allowing us to use the ICP-RIE machine. We thank F. Beil, J.P. Kotthaus, and F. Troiani for discussions.
warning/0506/astro-ph0506594.html	ar5iv	text	# The intermediate-age open clusters Ruprecht 61, Czernik 32, NGC 2225 and NGC 2262 ## 1 Introduction This paper continues a series dedicated to the study of open clusters in the third Galactic Quadrant, and aims at addressing fundamental questions like the structure of the spiral arms in this quadrant, and the precise definition of the Galactic disk radial abundance gradient outside the solar circle. A more detailed illustration of the motivations of this project are given in Moitinho (2001) and Baume et al (2004). Here we concentrate on four intermediate-age clusters (about the age of the Hyades - 600 Myrs - or older) Ruprecht 61 (VdB-Hagen 32), Czernik 32 (VdB-Hagen 11), NGC 2225 and NGC 2262, for which we provide the first photometric data and try to clarify their nature and to derive estimates of their fundamental parameters. In Table 1 we list the cluster coordinates, which we redetermined from Digital Sky Survey (DSS) maps on a visual inspection basis. The layout of the paper is as follows. Sect. 2 illustrates the observation and reduction strategies. An analysis of the geometrical structure and star counts in the field of the clusters is presented in Sect. 3, whereas a discussion of the Color-Magnitude Diagrams (CMDs) is performed in Sect. 4. Sect. 5 deals with the determination of clusters’ reddening, distance, metallicity and age and, finally, Sect. 6 summarizes our findings. ## 2 Observations and Data Reduction CCD $`BVI`$ observations were carried out with the CCD camera on-board the 1. 0m telescope at Cerro Tololo Interamerican Observatory (CTIO,Chile), on the nights of December 13 and 15, 2004. With a pixel size of $`0^{\prime \prime }.469`$, and a CCD size of 512 $`\times `$ 512 pixels, each pointing samples a $`4^{}.1\times 4^{}.1`$ field on the sky. The details of the observations are listed in Table 2 where the observed fields are reported together with the exposure times, the average seeing values and the range of air-masses during the observations. Figs. 1 to 4 show I=600 secs images obtained in the area of Ruprecht 61, Czernik 32, NGC 2225 and NGC 2262, respectively. Together with the clusters, we observed three control fields 15 arcmins apart from the nominal cluster centers to deal with field star contamination. Exposure of 600 secs in V and I were secured for these fields. The data have been reduced with the IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by NOAO, which are operated by AURA under cooperative agreement with the NSF. packages CCDRED, DAOPHOT, ALLSTAR and PHOTCAL using the point spread function (PSF) method (Stetson 1987). The two nights turned out to be photometric and very stable, and therefore we derived calibration equations for all the 130 standard stars observed during the two nights in the Landolt (1992) fields SA 95-41, PG 0231+051, Rubin 149, Rubin 152, T phe and SA 98-670 (see Table 2 for details). The adopted calibration equations are of the form: $`b=B+b_1+b_2\times X+b_3\times (BV)`$ $`v=V+v_1+v_2\times X+v_3\times (BV)`$ $`v=V+v_{1,i}+v_{2,i}\times X+v_{3,i}\times (VI)`$ $`i=I+i_1+i_2\times X+i_3\times (VI)`$ , where $`BVI`$ are standard magnitudes, $`bvi`$ are the instrumental ones and $`X`$ is the airmass; all the coefficient values are reported in Table 3. The standard stars in these fields provide the color coverage $`0.6(BV)2.2)`$. The final r.m.s. of the calibration are 0.031, 0.024 and 0.023 for the B, V and I filters, respectively. We generally used the third equation to calibrate the $`V`$ magnitude in order to get the same magnitude depth both in the cluster and in the field. Photometric errors have been estimated following closely Patat & Carraro (2001, Appendix A), which the reader is referred to for all the details. It turns out that the global photometric errors amount to 0.03, 0.05 and 0.20 at V = 12, 16 and 21.5 mag, respectively. The final photometric catalogs for Ruprecht 61, Czernik 32, NGC 2225 and NGC 2262 (coordinates, B, V and I magnitudes and errors) consist of 1166, 1047, 869 and 925 stars, respectively, and are made available in electronic form at the WEBDA<sup>2</sup><sup>2</sup>2http://obswww.unige.ch/webda/navigation.html site maintained by J.-C. Mermilliod. ## 3 Star counts and clusters’ sizes As we will show in this Section, our photometry covers entirely each cluster’s area allowing us to perform star counts to obtain improved estimates of the clusters’ size. In fact these clusters are generally very faint, poorly populated and compact (see Fig .1 to 4) and therefore could well fit within the CCD area. We derived the surface stellar density by performing star counts in concentric rings around the clusters’ nominal centers (see Table 1) and then dividing by their respective areas. Poisson errors have also been derived and normalized to the corresponding area. The field star contribution has been derived from the control field which we secured for each cluster, and the errors have been computed in the same way as for the cluster field. Ruprecht 61 The radial density profile for Ruprecht 61 is shown in Fig. 5 as a function of the V magnitude. Clearly, the cluster does not appear very concentrated, and it is deficient of bright stars. The cluster seems to emerge from the background in the magnitude range $`16V22`$. In this magnitude range the radius is not larger than 1.0 arcmin. We shall adopt the value of 1.0 arcmin as the radius of Ruprecht 61 throughout this paper. This estimate is in perfect agreement the value of 2.0 arcmin reported by Dias et al. (2002) for the cluster diameter. Czernik 32 The radial density profile for Czernik 32 is shown in Fig. 6 as a function of the V magnitude. The cluster shows a clear lack of bright stars ant it is mostly populated by stars of magnitude in the range $`18V22`$. In this magnitude range the radius is not larger than 1.0 arcmin. In conclusion, we are going to adopt the value of 1.0 arcmin as the Czernik 32 radius throughout this paper. This estimate is in reasonable agreement with the value of 3.0 arcmin reported by Dias et al. (2002) for the cluster diameter. NGC 2225 The radial density profile for NGC 2225 is shown in Fig. 7 as a function of the V magnitude. Also this cluster exhibits a deficiency of bright stars ant it is mostly populated by stars of magnitude in the range $`16V22`$. In this magnitude range the radius is not larger than 1.2 arcmin. This estimate is in reasonable agreement with (a bit smaller than) the value of 4.0 arcmin reported by Dias et al. (2002) for the cluster diameter. NGC 2262 The radial density profile for NGC 2262 is shown in Fig. 8 as a function of the V magnitude. This cluster is composed mainly by stars of magnitude in the range $`16V22`$, where the radius is around 1.0 arcmin. In conclusion, we are going to adopt the value of 1.0 arcmin as the radius of NGC 2262 throughout this paper. This estimate is much smaller than the value of 5.0 arcmin reported by Dias et al. (2002) for the cluster diameter. The estimates we provide for the radius, although reasonable, must be taken as preliminary. In fact, the size of the CCD is probably too small to derive conclusive estimates of the cluster sizes, that due to dynamical evolution and mass segregation tend to be normally under-estimated. Larger and deep field coverage is necessary to derive firmer estimates of the clusters radii. ## 4 The Colour-Magnitude Diagrams In Fig. 9 we present the CMDs obtained for the observed fields of the four clusters under investigation. All the stars observed in each field have been plotted (not only those within the derived cluster radii). In this figure, the open cluster Ruprecht 61 is shown together with the corresponding control field in the lower panels, Czernik 32 and NGC 2225 are presented in the middle panels, while finally NGC 2262 is presented in the upper panels. The control fields help us to better interpret these CMDs, which are clearly affected by strong foreground star contamination. Ruprecht 61. This cluster is presented in the lower panels of Fig. 9. It exhibits a Main Sequence (MS) extending from V=15-15.5, where the Turn Off Point (TO) is located, down to V=21.5. This MS is significantly wide, wider than photometric error at a given magnitude (see Sect. 2). We ascribe this to field star contamination, and to the presence of a sizeable binary star population, which mainly enlarge the MS toward red colors. However, the reality of this cluster seems to be secured by the shape of the MS with respect to the control field MS, whose population sharply decreases at V = 17. Also, the cluster MS is significantly bluer and more tilted than the field MS, which derives from the superposition of stars of different reddening located at all distances between the cluster and the Sun. Another interesting evidence is the possible presence of a clump of stars at V=14.5, which does not have a clear counterpart in the field, and which implies a cluster of intermediate-age. In fact if we use the age calibration from Carraro & Chiosi (2004), for a $`\mathrm{\Delta }V`$ (the magnitude difference between the red clump and the TO) of 0.5 mag, we infer an age around 1 Gyr. This estimate does not take into account the cluster metallicity, and therefore is simply a guess. In the following Sect. we shall provide a more robust estimate of the age through a detailed comparison with theoretical isochrones. Czernik 32. The open cluster Czernik 32 is presented in the lower-mid panels of Fig. 9. The TO located at $`V`$ 17, and a prominent clump at $`V`$ 16 with no counterpart in the field CMD are readily seen, yielding an estimated age of around 1.0 Gyr. The overall morphology of the CMDs is in this case very different from the field CMD leaving no doubt that Czernik 32 is a bona-fide intermediate-age open cluster. NGC 2225. The open cluster NGC 2225 is presented in the upper-mid panels of Fig. 9. Again, the overall morphology of the cluster CMDs, with a TO and an evident red clump, is very different from the field CMD. Indeed, the field sequence is much less populated and stops at V $``$ 16.5, giving a first impression that this is a bona-fide intermediate-age open cluster. The TO located at V $``$ 16, and red clump at V $``$ 15, allow to estimate an age of around 1.0 billion years, confirming the first impression that NGC 2225 is an intermediate age cluster. NGC 2262. The open cluster NGC 2262 is presented in the upper panels of Fig. 9. The cluster’s CMD reveals a TO is located at $`V`$ 16, and a possible clump at $`V`$ 16 as well, which provide a rough estimate of around 0.5 Gyrs for the age of NGC 2262. The overall morphology of the CMDs is also in this case very different from the field CMD confirming that this is a bona-fide intermediate-age open cluster. ## 5 Deriving clusters’ fundamental parameters In this section we are going to perform a detailed comparison of the star distribution in the clusters’ CMDs with theoretical isochrones. For this study, we adopt in this study the Padova library from Girardi et al. (2000). This comparison is clearly not an easy exercise. In fact, the detailed shape and position of the various features in the CMD (MS, TO and clump basically) depends mostly on age and metallicity, and then also on reddening and distance. The complex interplay between the various parameters is however well known, and we refer to Chiosi et al. (1992) and Carraro (2005) as nice examples of the underlying technique. Our basic strategy is to survey different age and metallicity isochrones in an attempt to provide the best fit of all the CMD features both in the $`V`$ vs $`(BV)`$ and in the $`V`$ vs $`(VI)`$ CMD. To further facilitate the fitting procedure, by increasing the contrast between the cluster and the field population, we shall consider only the stars which lie within the cluster radius as derived in Sect. 3. Finally, to derive the clusters’ distances from reddening and apparent distance modulus, a reddening law must be specified. In this study we shall adopt the normal reddening law $`Av=3.1\times E(BV)`$ in deriving the clusters’ distances. Additionally to finding the best fit, we also estimated the uncertainties in the basic parameters. These uncertainties simply reflect the range of parameters that yields a reasonable fit to the clusters CMDs. The errors are reported in Table 4, and an example of the procedure is shown in Fig.10 for the case of Ruprecht 61. The best fir for all the clusters, achieved simultaneously in the V vs $`(BV)`$ and in the $`V`$ vs $`(VI)`$ planes, are shown in Figs. 11-14. Ruprecht 61. The fitting procedure and the isochrone solution for this cluster are shown in Figs. 10 and 11. We obtained the best fit for an age of 1.3 Gyrs and a metallicity Z=0.008 (see middle panel of Fig. 11, and Fig. 12). In fact, the shape of the TO in the left panel of Fig. 11 (for the Z=0.004 isochrone) is clearly different from the underlying cluster sequence, and the same can be noticed for the Z=0.019 isochrone (right panel), where the red hook shows a shape which does not fit very well the star distribution in the cluster. In details, the red hook is too red and somewhat faint with respect to the Z=0.008 isochrone and the actual stars distribution. To get a bluer and brighter red hook, one should use a younger isochrone, which will however possess a too red clump and RGB, when fixed to the TO. The inferred reddening and apparent distance modulus are E(B-V)=0.30 (E(V-I)=0.41, right panel in Fig. 12) and (m-M)=13.85, respectively. As a consequence, the cluster possesses a heliocentric distance of 3.9 kpc, and is located at a Galactocentric distance of 9.4 kpc, assuming 8.5 kpc as the distance of the Sun to the Galactic Center. Czernik 32. The isochrone solution for this cluster is displayed in Fig.12. We obtained the best fit for an age of 1 Gyr and a metallicity Z=0.008. The inferred reddening and apparent distance modulus are E(B-V)=0.85 (E(V-I)=1.08) and (m-M)=15.7, respectively. These values situate the cluster at a heliocentric distance of 4.1 kpc, which corresponds to a Galactocentric distance of 10.8 kpc. The overall fit is also good in this case, the detailed shape of the MS and TO are nicely reproduced, as well as the color of the clump. NGC 2225. The isochrone solution for this cluster is presented in Fig.13. We obtained the best fit for an age of 1 Gyr and a metallicity Z=0.008, which reproduces the sharp cluster sequence extremely well. The inferred reddening and apparent distance modulus are E(B-V)=0.35 (E(V-I)=0.50) and (m-M)=13.6, respectively. Therefore the cluster has a heliocentric distance of 3.2 kpc, and is located at a Galactocentric distance of 11.2 kpc. NGC 2262. The isochrone solution for this cluster is shown in Fig.14. We obtained the best fit for an age of 1 Gyr and a metallicity of Z=0.008. The inferred reddening and apparent distance modulus are E(B-V)=0.55 (E(V-I)=0.72) and (m-M)=14.5, respectively, which put the cluster at a heliocentric distance of 3.6 kpc, or at a Galactocentric distance of 11.7 kpc. ## 6 Conclusions We have presented the first CCD $`BVI`$ photometric study of the star clusters Ruprecht 61, Czernik 32, NGC 2225 and NGC 2262. Through a star count analysis we have refined previous estimates of the clusters’ coordinates and apparent radii. A detailed comparison of the clusters’ CMDs with theoretical isochrones has allowed us to infer the aggregates’ basic parameters and their uncertainties , which are summarized in Table 4. In detail, the fundamental findings of this paper are: $``$ the best fit reddening estimates support within the errors a normal extinction law toward the four clusters; $``$ all the clusters turn out to be of intermediate age, and not far from the Sun toward the anti-center direction. $``$ the photometric estimates of the metallicity are lower than solar, as expected for clusters located between 9 and 11 kpc from the Galactic Center (Carraro et al. 1998). ## Acknowledgements The observations presented in this paper have been carried out at Cerro Tololo Interamerican Observatory CTIO (Chile). CTIO is operated by the Association of Universities for Research in Astronomy, Inc. (AURA), under a cooperative agreement with the National Science Foundation as part of the National Optical Astronomy Observatory (NOAO). The work of G.C. is supported by Fundación Andes. D.G. gratefully acknowledges support from the Chilean Centro de Astrofísica FONDAP No. 15010003. This work has been also developed in the framework of the Programa Científico-Tecnológico Argentino-Italiano SECYT-MAE Código: IT/PA03 - UIII/077 - período 2004-2005. A.M. acknowledges support from FCT (Portugal) through grant SFRH/BPD/19105/2004. This study made use of Simbad and WEBDA databases.
warning/0506/quant-ph0506105.html	ar5iv	text	# Modified Grover’s search algorithm for the cases where the number of solutions is known ## 1 Introduction With the advent of quantum computation many quantum algorithms \[References-References\] which, work faster than their classical counter parts, have appeared. Among these quantum algorithms, Grover’s algorithm \[References\] got special attention of the whole community because of its wide applicability in searching databases. Actually, searching databases is one of the most important problems in computer science and real life. This fact has motivated people to develop a large number of algorithms to search different kind of databases \[References\]. We are interested about a database of $`N`$ unsorted items, having $`M`$ solutions (where $`MN)`$. Any classical algorithm takes $`O(N)`$ steps to search such a database. Grover’s quantum algorithm \[References\] searches such a database in $`O(\sqrt{N/M})`$ steps. Till now, the time complexity of Grover’s algorithm is minimum among all the algorithms designed for the same purpose. Tight bound on Grover searching has been studied by many people \[References\] but it does not establish any tight bound on quantum search in general. This fact has motivated us to explore the possibility of improvement in special cases. In Grover’s original algorithm the number of solution $`M`$ is unknown. In the present work we consider a special case of Grover’s algorithm and assume that either $`M`$ is known or an upper bound of $`M`$ is known. In the first case time complexity reduces to $`O(logN)`$ and in the second case it reduces to $`O(MlogN)`$. The reduction is considerably large when $`M`$ is small and $`N`$ is large. In section 2 we briefly discuss Grover’s algorithm. The modified algorithm, is discussed in section 3. In section 4 we have discussed time complexity of various cases. Finally we conclude in section 5. ## 2 Grover’s algorithm As we have already stated, we are interested to search a database of $`N`$ items out of which $`M`$ are the solutions. In Grover’s search algorithm we assign an index to each element and search on those indices. Now for our convenience if we assume that $`N=2^n`$ then we can store all the indices in $`n`$ qubits since the indices varies from $`0`$ to $`N1`$. A particular instance of the search problem can conveniently be represented by a function $`f`$, which takes an integer $`x`$, in the range $`0`$ to $`N1`$. By definition, $`f(x)=1`$ if $`x`$ is a solution to the search problem and $`f(x)=0`$ if $`x`$ is not a solution to the search problem. Grover’s algorithm uses an unitary operator as a quantum oracle which flips the oracle qubit if $`f(x)=1`$. Essentially, the Oracle marks the solutions to the search problem by shifting the phase of the solution. The search oracle is applied only $`O(\sqrt{N/M})`$ times in order to obtain a solution on a quantum computer. This is done in following steps: 1. The algorithm begins with a quantum register in the state $`\|0^n`$. 2. The Hadamard transform is used to put the register in a equal superposition of $`N=2^n`$ states. This is how we used to prepare the input state $`\|x>`$ for the oracle. 3. A quantum subroutine, known as the Grover iteration is repeatedly applied. The Grover iteration may be broken in following four steps: a) Apply the oracle b) Apply the Hadamard transformation on $`n`$ qubits c) Perform a conditional phase shift on the computer, with every computational basis state except $`\|0`$ receiving a phase shift of $`1`$ d) Apply the Hadamard transformation on $`n`$ qubits. ## 3 The modified algorithm To simplify the understanding of Grover’s algorithm, we can assume that the initial superposition is constituted of 2 parts: the solution states and the non-solution states and represent the state as $$\|S=\mathrm{cos}\theta \|0>+\mathrm{sin}\theta \|1$$ (1) where $`\|0`$ represents the non-solution states and $`\|1`$ represents the solution states and $$\mathrm{cos}\theta =\sqrt{\frac{(NM)}{N}}$$ $$\mathrm{sin}\theta =\sqrt{\frac{M}{N}}.$$ The Oracle can be considered as a $`2\times 2`$ matrix which flips the phase of the solution states. It can be written as : $$O=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$ At the end of each Grover iteration (Grover iteration is a phase flip of the solution states, followed by an inversion of all states about the mean), the initial state gets rotated by an angle of $`2\theta `$ in a direction such that it moves closer to the solution states. In other words, each Grover iteration increases the probability of the solution states (simultaneously decreasing the probability of the non-solution states). Therefore, the correct solution can be measured with a high probability after a certain number of Grover iterations. Essentially, a particular Grover iterator redistributes the probability among the possible states in two steps. First it flips the phase of the solution states and then inverts about the mean. In this process, the probability of nonsolution states gets reduced and the reduced probability is added to those of the solution states. Here an important question arises: Is it essential to invert the states about mean? The answer is no! Actually, the essential condition is unitarity of the operation. When $`M`$ is unknown, then this one of the unitary operation through which we can redistribute probability according to the requirement and conserve the total probability. So when $`M`$ is unknown the state represented by (1) has to be rotated by $`2\theta `$ in each step. But if we know the, value of $`M`$ (i.e. we know $`\theta `$), then we can introduce an unitary operation which vanishes the probability of appearance of nonsolution states and uniformly distributes that probability among all the solution states. This new unitary operation exploits the fact that if $`M`$ is known then the amount of rotation which can map the initial state into the solution state is known. Geometrically, these means an inversion about a suitable point (instead of the inversion about the mean). The equation to determine the number of iterations $`I`$ required in conventional Grover’s algorithm is $$\theta +I(2\theta )=\frac{\pi }{2}.$$ (2) Now, instead of carrying out Grover’s iteration large number of times, we propose carrying out the same action in one step i.e. instead of rotating the current search state by $`2\theta `$, we propose rotating it directly by $`k\theta `$ where: $$\theta +k\theta =\frac{\pi }{2}$$ (3) i.e. $$k\theta =\frac{\pi }{2}\theta .$$ (4) Thus, if we can rotate the current search state (initial state) by $`k\theta `$ then we can obtain the desired solution state in a single iteration. The time complexity of the process is $`O(logN)`$ (to create Hadamard superposition). A $`2\times 2`$ matrix that rotates a state vector (represented by a $`2\times 1`$ matrix in 2 dimensions) by $`k\theta `$ can be written as: $$\left(\begin{array}{cc}cosk\theta & sink\theta \\ sink\theta & cosk\theta \end{array}\right).$$ Replacing $`k\theta `$ by $`\frac{\pi }{2}\theta `$ we get a new operator $`A`$ defined as follows: $$A=\left(\begin{array}{cc}sin\theta & cos\theta \\ cos\theta & sin\theta \end{array}\right).$$ (5) To understand the physical meaning of this operation let us assume that the rotation operation A is obtained as the oracle operation O followed by another operation (say X), i.e. XO=A. Solving the above equation, we get X as: $$X=\left(\begin{array}{cc}sin\theta & cos\theta \\ cos\theta & sin\theta \end{array}\right).$$ (6) The matrix X can be written in operator form as: $$X:=\left(sin\theta \|0+cos\theta \|1\right)0\|+\left(cos\theta \|0sin\theta \|1\right)1\|.$$ (7) According to our basic assumption the value of $`M`$ and $`N`$ are known. Therefore, $`cos\theta `$ and $`sin\theta `$ are known and we can prepare the unitary operation $`X`$. It is easy to check that $`X`$ is unitary and physically $`X:`$ represents a quantum gate which causes an inversion about a point such that the nonsolution state probabilities are reduced to zero. The $`X`$ can be multiplied with $`O`$ to produce $`A`$, which operates on $`\|S`$as follows, $$A:\|S=A:\left(\begin{array}{c}cos\theta \\ sin\theta \end{array}\right)=\left(\begin{array}{c}0\\ 1\end{array}\right).$$ (8) Thus, we are left only with solution states that can be obtained by performing a measurement. Essentially, the modification of the point of inversion reduces the time complexity in our case. But only if we know the total number of solutions then we can choose the suitable point of inversion. ## 4 Time complexity in various cases: 1. M (the number of solutions to the search query) is known Only 1 iteration is required to reach the solution state. The input is prepared by applying appropriate number of Hadamard gates resulting in equal superposition of N states. Thus, the total time complexity would be $`O(logN)`$. 2. M is unknown but we can estimate an upper bound on the possible value of M. The algorithm suggested above can only be executed only for a particular value of M. Since we are aware of the upper bound, measurement of the register which holds the answer (for a particular value of M) is checked to be correct. Thus, for each value of M we are required to verify the correctness of the answer provided by running the algorithm. The answers obtained can be checked easily as stated in \[References\]. This approach will lead to a time complexity of $`O(MlogN)`$. 3. M is unknown and we cannot estimate an upper bound on the possible value of M. An alternative approach can be adopted in this case. Recently we have given a proposal \[References\] to handle this case using concurrency control techniques and marking. This proposal reduces the complexity to $`O(M+logN)`$. ## 5 Conclusion This paper proposes a scheme to search a database of $`N`$ unordered items in $`O(logN)`$ when $`M`$ is known. And in $`O(MlogN)`$ when $`M`$ is unknown but an estimation of upper limit of $`M`$ is possible. This improvement in complexity is considerable and it will be more prominent with the increase of the size of the database. There exist many applications of Grover’s algorithm. Thus, an improvement in Grover’s search will result in the improvement in time complexity of all these applications. This is a special case of Grover’s algorithm where time complexity is less than that of the conventional Grover’s algorithm. There may exist many similar special cases of more general quantum search problem where complexity is less. Thus the present study opens up a possibility to look at the quantum search problems from a new perspective.
warning/0506/cond-mat0506507.html	ar5iv	text	# Stabilization of nonlinear velocity profiles in athermal systems undergoing planar shear flow ## I Introduction The response of fluids to applied shear is an important and well studied problem. When a Newtonian fluid is slowly sheared, for example by moving the top boundary of the system at fixed velocity $`u`$ in the $`x`$-direction at height $`y=L_y`$ relative to a stationary bottom boundary at $`y=0`$, a linear velocity profile $`v_x(y)=\dot{\gamma }y`$ is established, where the shear rate is $`\dot{\gamma }=u/L_y`$. In addition, in the small shear rate limit the shear stress is linearly related to the shear rate $`\sigma _{xy}=\eta \dot{\gamma }`$, where $`\eta `$ is the shear viscosity. This relation is often employed to measure the shear viscosity of simple fluids. However, what is the response of non-Newtonian fluids like granular materials to an applied shear? In contrast to simple liquids, it is extremely difficult to predict the response of dense granular media to shear because they are inherently out of thermal equilibrium, interact via frictional and enduring contacts, and possess a nonzero yield stress. Granular materials do not flow homogeneously when they are sheared, instead, shear is often localized into shear bands. When this occurs, most of the flow is confined to a narrow, locally dilated region near the shearing boundary while the remainder of the system is nearly static. Recent experimental work investigating the response of dense granular materials to shear includes studies of Couette flow in 2D howell , in 3D for spherical particles losert ; mueth2 and as a function of particle shape mueth , studies of cyclic planar shear mueggenburg , studies of wide shear zones formed in the bulk using a modified Couette geometry fenistein , and studies of chute flow pouliquen . Recent theoretical studies jenkins have shown that kinetic theory can correctly predict the velocity, granular temperature, and density profiles found in experiments of Couette flow in the dilute regime wildman . However, kinetic theory and modifications included to account for diverging viscosity are not analytically tractable, and are unlikely to predict accurately the properties of dense shear flows. Thus, molecular dynamics simulations are often employed to study dense granular shear flows, for example in Refs. thompson ; aharonov ; volfson ; dacruz . We choose a similar plan of attack and perform molecular dynamics simulations of model frictionless dense granular systems undergoing boundary-driven planar shear flow in 2D. We are interested in answering several important questions: What are the velocity, granular temperature, and density profiles as a function of the velocity of the shearing boundary? In particular, do highly localized velocity profiles form and, if so, are they stable at long times? We will investigate these questions in simple 2D systems composed of inelastic but frictionless particles in planar shear cells, and in the absence of gravity. Our intent is to understand in detail the time and spatial dependence of the response in these simple systems to shear first, and then extend our studies to 3D systems, systems composed of frictional particles, and Couette shear cells. In Sec. VI, we present preliminary results from these future studies. We previously reported that nonlinear velocity profiles form in repulsive athermal systems when the velocity of the shearing wall exceeds the speed of shear waves in the material xu . In the present article, we study much longer time scales and show that nonlinear velocity profiles are unstable at long times. In addition, the granular temperature and density profiles become uniform throughout the system at long times. Thus, granular temperature and density gradients cannot be maintained by planar shear flow in dense systems with dissipative but frictionless interactions. We measure the time $`t_l`$ for the velocity profiles to become linear and find that $`t_l`$ scales as a power-law in the velocity of the shearing boundary, and increases rapidly as the average density of the system approaches random close packing $`\varphi _{\mathrm{rcp}}`$ from above, where $`\varphi _{\mathrm{rcp}}0.84`$ in 2D ohern . To maintain a granular temperature difference across the system, we also studied systems that were both sheared and vertically vibrated. We find that if the difference in the granular temperature across the system exceeds a threshold that is comparable to the glass transition temperature in an equilibrium system at the same average density, nonlinear velocity profiles are stable at long times. Finally, we show that shear bands, or highly localized velocity profiles, form in the vibrated and sheared systems if the shear stress is tuned below the yield stress of the static part of the system. ## II Methods Before discussing our results further, we will first describe our numerical model and methods. We performed molecular dynamics simulations of purely repulsive and frictionless athermal systems undergoing boundary-driven planar shear flow in two spatial dimensions. The systems were composed of $`N/2`$ large and $`N/2`$ small particles with equal mass $`m`$ and diameter ratio $`1.4`$ to prevent crystallization and segregation during shear. The starting configurations were prepared by choosing an average packing fraction and random initial positions and then allowing the system to relax to the nearest local potential energy minimum ohern using the conjugate gradient method numrec . During the quench, periodic boundary conditions were implemented in both the $`x`$\- and $`y`$-directions. Following the quench, particles with $`y`$-coordinates $`y>L_y`$ ($`y<0`$) were chosen to comprise the top (bottom) boundary. Thus, the top and bottom walls were rough and amorphous. Shear flow in the $`x`$-direction with a shear gradient in the $`y`$-direction and global shear rate $`u/L_y`$ was created by moving all particles in the top wall at fixed velocity $`u`$ in the $`x`$-direction relative to the stationary bottom wall. During shear flow, periodic boundary conditions were imposed in the $`x`$-direction. We chose an aspect ratio $`L_x/L_y=1/4`$ with more than $`50`$ particles along the shear gradient direction to reduce finite-size effects, and focused on systems with packing fractions in the range $`\varphi =[0.835,0.95]`$. Both bulk particles and particles comprising the boundary interact via the purely repulsive harmonic spring potential $$V(r_{ij})=\frac{ϵ}{2}\left(1\frac{r_{ij}}{\sigma _{ij}}\right)^2\mathrm{\Theta }\left(\frac{\sigma _{ij}}{r_{ij}}1\right),$$ (1) where $`ϵ`$ is the characteristic energy scale of the interaction, $`\sigma _{ij}=(\sigma _i+\sigma _j)/2`$ is the average diameter of particles $`i`$ and $`j`$, $`r_{ij}`$ is their separation, and $`\mathrm{\Theta }(x)`$ is the Heaviside step function. Note that the interaction potential is zero when $`r_{ij}\sigma _{ij}`$. In our simulations, we employ athermal or dissipative dynamics with no frictional or tangential forces herrmann . The position and velocity of particles in the bulk were obtained by solving $$m\frac{d^2\stackrel{}{r}_i}{dt^2}=\stackrel{}{F}_i^rb_n\underset{j}{}\left[\left(\stackrel{}{v}_i\stackrel{}{v}_j\right)\widehat{r}_{ij}\right]\widehat{r}_{ij},$$ (2) where $`\stackrel{}{F}_i^r=_jdV(r_{ij})/dr_{ij}\widehat{r}_{ij}`$, the sums over $`j`$ only include particles that overlap $`i`$, $`\stackrel{}{v}_i`$ is the velocity of particle $`i`$, and $`b_n>0`$ is the damping coefficient. We focus on underdamped systems in this study; specifically, for most simulations we use $`b_n=0.0375`$ (coefficient of restitution $`e=0.92`$). In addition, we show some results that cover a range of $`e=[0.1,0.99]`$. The units of length, energy, and time are chosen as the small particle diameter $`\sigma `$, $`ϵ`$, and $`1/\omega _c\sigma \sqrt{m/ϵ}`$, and all quantities were normalized by these. We also performed simulations in which the system was both sheared and vertically vibrated. As discussed above, the system was sheared by constraining the top boundary to move in the $`x`$-direction at fixed speed $`u`$ relative to the bottom boundary. In addition, particles in either the top or bottom boundary were vibrated vertically such that the $`y`$-coordinates of the boundary particles varied in time as $$y_i=y_{i0}+A\mathrm{sin}\left(\omega t\right),$$ (3) where $`y_{i0}`$ is the initial position of boundary particle $`i`$, $`A`$ is the amplitude, and $`\omega `$ is the angular frequency of the vibration. We fixed $`A=\sigma /5`$ so that the vibrations do not lead to large, unphysical particle overlaps and $`\omega `$ was tuned over a range of frequencies. Note that in our choice of units $`\omega `$ is normalized by the natural frequency $`\omega _c`$ of the linear spring interactions. It should also be pointed out that the simulations with both shear and vibration are not performed at constant volume but at fixed average volume. We measured several physical quantities in these simulations including the local flow velocity $`v_x`$, packing fraction $`\varphi `$, and velocity fluctuations or granular temperature $`(\delta v_y)^2=v_y^2v_y^2`$ in the shear gradient direction, as a function of the boundary velocity $`u`$ and vibration frequency $`\omega `$. To study properties as a function of the vertical distance from the fixed boundaries, we divided the system into rectangular bins centered at height $`y`$ and averaged the quantities over the height of the bin $`\mathrm{\Delta }y2`$ large particle diameters. To improve the statistics and to study time dependence, we also performed ensemble averages over at least $`50`$ different initial configurations. In the discussion below, we will denote ensemble averages using $`.`$. When the systems are stationary but still fluctuate in time, we also perform time averages and these are denoted by $`._t`$. We also monitored the shear stress $`\sigma _{xy}`$ during the course of our simulations. The total shear stress in the bulk can be calculated using the virial expression $$\sigma _{xy}=\frac{1}{L_xL_y}\left(\underset{i}{}\delta v_{i,x}\delta v_{i,y}+\underset{i>j}{}r_{ij,x}F_{ij,y}\right),$$ (4) where $`r_{ij,x}`$ is the $`x`$-component of $`\stackrel{}{r}_{ij}=\stackrel{}{r}_i\stackrel{}{r}_j`$, $`F_{ij,y}`$ is the $`y`$-component of the total pair force $`\stackrel{}{F}_{ij}`$ including both conservative and damping forces, and $`\delta v_{i,x}`$ and $`\delta v_{i,y}`$ measure deviations in the velocity of a particle from the velocity averaged over the entire system. We also compared the shear stress obtained from Eq. 4 with the shear stress $`\|F_x^B\|/L_x`$ on the boundaries, where $`F_x^B`$ is the total force in the $`x`$-direction acting on either the top or bottom boundary. On average, the bulk shear stress in Eq. 4 and the shear stress on the boundaries were within a few percent of each other. To make contact with the recent results on sheared Lennard-Jones glasses in Ref. varnik , we measured the yield stress of the static starting configurations. To do this, we applied a constant horizontal force $`F`$ (or shear stress $`\sigma _{xy}=F/L_x`$) to the top boundary. Results from the simulations at constant horizontal force are shown in Fig. 1. Initially, the system flows. However, when the shear stress is below the yield stress, the system finds a state that can sustain the applied shear stress and stops flowing. If the applied shear stress is above the yield stress, the system will flow indefinitely. We defined the yield shear stress $`\sigma _0`$ as the minimum shear stress above which the shear strain continues to increase beyond $`\gamma =10`$. ## III Results In this section, we report the results from two sets of numerical simulations of purely repulsive and frictionless athermal systems. Section III.1 presents results from simulations of boundary-driven planar shear flow. We show that nonlinear velocity profiles form at short time scales, but they slowly evolve into linear profiles at long times. As the velocity profiles evolve toward linear ones, the local granular temperature and packing fraction become uniform. Section III.2 presents results from simulations of boundary-driven planar shear flow in the presence of vertical vibrations. We find that highly nonlinear velocity profiles can be stabilized at long times if a sufficiently large granular temperature difference is maintained across the system. ### III.1 Time evolution of velocity profiles In our previous studies of boundary-driven planar shear flow xu , we reported that nonlinear velocity profiles form when the velocity $`u`$ of the shearing boundary exceeds $`u_c=u_s/2`$, where $`u_s`$ is speed of shear waves in the system. This condition was obtained by comparing the time $`t_s=2L_y/u_s`$ for a shear wave to traverse the system to the time $`t_u=L_y/u`$ for the system to shear unit strain. A rough estimate of the speed of shear waves (at least in the low shear rate limit) can be obtained from $`\sqrt{G/\rho }`$, where $`G`$ is the static shear modulus and $`\rho `$ is the mass density. A more precise way to measure $`u_s`$ is to calculate the transverse current correlation function $`C_T(k,\nu )`$ remark4 as a function of frequency $`\nu `$ and wave number $`k=2\pi n/L_x`$ ($`n=`$ integer) and determine the slope of the resulting dispersion relation $`\nu (k)`$ hansen . One of the novel aspects of our previous work was that we showed that underdamped systems would be extremely susceptible to nonlinear velocity profiles near random close packing since the critical velocity $`u_c`$ tends to zero at $`\varphi _{\mathrm{rcp}}`$. We therefore predicted that any $`u>0`$ would give rise to nonlinear velocity profiles in these systems near $`\varphi _{\mathrm{rcp}}`$. In these prior studies, we measured the velocity, packing fraction, and mean-square velocity fluctuation profiles after shearing the system for a strain of at least $`5`$ and times $`t>t_s`$ so that the shear stress had relaxed to its long-time average value. This protocol for bringing sheared systems to steady-state is typical in both simulations and experiments. If $`t_s`$ were the only relevant time scale, the nonlinear velocity profiles that occur in boundary-driven planar shear flow when $`u>u_c`$ would be stable over long times. However, in more recent studies, we have found that these nonlinear velocity profiles are not stable at long times $`tt_s`$ and slowly evolve toward linear profiles. In Fig. 2, we show the slow time evolution of the velocity profile in a system sheared at $`u=0.364`$, $`\varphi =0.90`$, and $`e=0.92`$. Note that the strains beyond which the profiles become linear are extremely large, $`\gamma >25`$. We will also show below that the time required for the velocity profile to become linear increases as the system becomes more elastic. Most of our previous work focused on nearly elastic systems with $`e0.98`$ and timescales near $`t_s`$, which made it difficult to detect the profile’s slow evolution. It would be interesting to know whether this slow evolution can be seen in experiments on sheared granular media, or other particulate systems. However, few experiments have studied such large strain and time scales. We note that recent experiments mueggenburg on granular media undergoing planar cyclic shear do show slow evolution; however, further experiments are required. Similarly, simulations of frictional granular media undergoing Couette flow in 3D have reported significant time evolution of the measured velocity profiles baran . To quantify the shape of the velocity profiles as a function of time, we define the degree of nonlinearity as $$\mathrm{\Delta }(t)=\left\|1\frac{v_x(y,t)}{v_x(y)_t}\right\|,$$ (5) where $`v_x(y)_t`$ is the time and ensemble average of the velocity in the flow direction after the evolution of the profile has ended. To measure the degree of nonlinearity, we calculated $`\mathrm{\Delta }(t)_y`$ averaged over the central part of the system excluding layers that are immediately adjacent to the top and bottom walls. $`\mathrm{\Delta }(t)_y0`$ corresponds to a nearly linear velocity profile, while $`\mathrm{\Delta }(t)_y1`$ is highly nonlinear. In Fig. 3, we show $`\mathrm{\Delta }(t)_y`$ for a system with $`\varphi =0.90`$ and $`e=0.92`$ at several different velocities of the shearing boundary. At each $`u`$, the velocity profiles slowly evolve toward linear profiles and $`\mathrm{\Delta }(t)_y`$ decays to zero at long times. To characterize the long-time behavior, we define a timescale $`t_l`$ as the time required for $`\mathrm{\Delta }(t)_y`$ to decay to $`0.05`$, which is slightly above the noise level. By definition, the velocity profiles are steady and linear for times $`t>t_l`$. At large $`u`$, the velocity profiles approach the linear profile from below as shown in Fig. 2. At small $`u`$, $`\mathrm{\Delta }(t)_y`$ decays to zero quickly, but the profiles jump above and below the linear profile as $`tt_l`$. The fact that the velocity profile has significant excursions above and below linearity at short times signals that the system is in the quasistatic flow regime. This behavior at small $`u`$ is interesting but not the topic of this study. In Fig. 4, we show that the time $`t_l`$ required for the velocity profile to become linear increases with the velocity $`u`$ of the shearing boundary. More precisely, $`t_lt_s`$ appears to scale as a power-law in $`uu_l`$ $$t_lt_s\left(uu_l\right)^\beta ,$$ (6) for $`u>u_l`$, where $`\beta 0.48\pm 0.02`$. For $`u<u_l`$, we find that the velocity profiles are nonlinear only for times $`t<t_s`$, and are linear for all subsequent times. We also investigated the sensitivity of $`t_l`$ to changes in the packing fraction and damping coefficient. In Fig. 5 (a), we show that $`t_l`$ increases sharply near random close packing $`\varphi _{\mathrm{rcp}}`$ at small damping. This rise in $`t_l`$ is suppressed at large damping coefficients as shown in Fig. 5. The rapid rise in $`t_l`$ at least for underdamped systems at densities below random close packing again suggests that these systems are susceptible to the formation of strongly nonlinear velocity profiles. These results are surprising and raise an important question regarding the physical mechanism that is responsible for the slow evolution of the velocity profiles. As a first step in addressing this question, we show below that granular temperature differences across the system give rise to nonlinear velocity profiles, and that, if a sufficiently large granular temperature difference can be maintained, nonlinear velocity profiles will be stable at long times. ### III.2 Combining vibration and shear: Stabilizing nonlinear velocity profiles at long times Recent experiments on sheared granular materials have shown that strongly nonlinear velocity profiles are accompanied by spatially dependent granular temperature profiles losert . Thus, an important question to ask is what role does the granular temperature play in determining the shape of the velocity profiles in sheared granular systems. Also, do these systems require a sufficiently large granular temperature difference across the system to possess strongly nonlinear velocity profiles? Our results in Fig. 6 suggest that the shapes of the granular temperature and velocity profiles are strongly linked. In this figure, we show the time evolution of the mean-square velocity fluctuations in the shear-gradient direction, $`\delta v_y^2`$, following the initiation of shear. The sequence of times is identical to that shown in Fig. 2. At short times, there is a large difference in the mean-square velocity fluctuations between the ‘hot’ shearing boundary and ‘cold’ stationary boundary, and the velocity profile is highly nonlinear. In contrast, at long times, the mean-square velocity fluctuations are uniform and the velocity profile is linear. Note that in the systems studied here, the granular temperature difference across the system, $`\mathrm{\Delta }T`$, is roughly equal to the granular temperature near top boundary, since the mean-square velocity fluctuations are much smaller near the bottom stationary boundary. Figs. 2 and 6 show that large granular temperature differences and nonlinear velocity profiles occur together. However, under steady shear, the granular temperature profiles become uniform and the velocity profiles become linear at long times. To further investigate the connection between the granular temperature and velocity profiles, we study systems that are both sheared and vibrated, and are thus designed so that granular temperature differences across the system can be maintained at long times. As discussed above in Sec. II, in our second set of simulations we drive the top boundary at constant horizontal velocity $`u`$, while the location of each particle in the top or bottom boundary oscillates sinusoidally in time with fixed amplitude $`A`$ and frequency $`\omega `$. The vertical vibrations cause the mean-square velocity fluctuations and packing fraction to become spatially nonuniform with higher velocity fluctuations and lower packing fraction or dilatancy near the vibrated boundary. Fig. 7 clearly demonstrates that vertical vibration coupled with shear flow gives rise to stable nonlinear velocity profiles at long times. At low and also at high vibration frequencies, the degree of nonlinearity $`\mathrm{\Delta }(t)_y`$ of the velocity profile decays to small values at long times, as we found previously in Fig. 3 for the unvibrated systems. In contrast, at frequencies near $`\omega _c`$, $`\mathrm{\Delta }(t)_y`$ is nonzero at long times. We note that the longest times shown in Fig. 7 are at least $`10`$ times longer than $`t_l`$ in the corresponding unvibrated system, confirming that these are indeed steady-state results. The results presented in Figs. 8 (a) and (b) suggest that large and sustained granular temperature differences across the system and the resulting dilatancy are responsible for stable nonlinear velocity profiles. This figure shows that nonlinear velocity profiles occur when there are large granular temperature differences at $`\omega =0.4`$, $`0.8`$, and $`1.4`$, while linear profiles are found when the velocity fluctuations are uniform at $`\omega =0.1`$. Moreover, the degree of nonlinearity in the velocity profiles increases with the magnitude of the granular temperature difference across the system. Fig. 8 (a) also shows the glass transition temperature at the same average packing fraction; we discuss the relevance of this temperature in more detail in Sec. IV remark . The granular temperature difference across the system increases with $`\omega `$ for $`\omega <\omega ^{}`$ but decreases when $`\omega >\omega ^{}`$ with $`\omega ^{}\omega _c`$. The largest granular temperature difference occurs near the natural frequency $`\omega _c`$ of the repulsive spring interactions remark2 . The decrease for $`\omega >\omega ^{}`$ occurs because particles adjacent to the vibrating boundary do not have enough time to react to the collision with the boundary before another collision occurs. As a result, vibrations at large frequencies simply reduce the effective height of the system by the amplitude of the vibration but do not induce large granular temperature differences. We note that $`\omega ^{}`$ does not appear to depend on the shearing velocity $`u`$. It is important to emphasize the fact that differences in the mean-square velocity fluctuations across the system, not the magnitude of the fluctuations themselves, are important in determining the shape of the velocity profiles. For example, a system with the same interactions, average density, and uniform temperature $`\delta v_y^2>10^3`$ will possess a linear velocity profile when sheared over the same range of $`u`$. Similarly, the system vibrated at $`\omega =0.1`$ (shown as circles in Fig. 8) is in a glassy state with small but relatively uniform mean-square velocity fluctuations $`\delta v_y^210^4`$ and also possesses a linear velocity profile. In systems that possess nonlinear velocity profiles, we find that the largest local shear rate does not occur equally likely at both boundaries. Instead, the portion of the system with the largest local shear rate always forms near the boundary with the largest $`\delta v_y^2`$ and resulting dilatancy. This is confirmed in Fig. 9, which shows results in a system that is identical to that in Fig. 8 except the bottom wall, not the top wall, is vibrated. The vertical vibrations induce a larger granular temperature near the bottom vibrated but unsheared boundary. Thus, the largest local shear rates occur near the bottom boundary, as shown for $`\omega =0.4`$, $`0.8`$, and $`1.4`$ in Fig. 9 (c). ## IV Discussion In this section, we will focus on several aspects of our results in more detail. First, we find that nonlinear velocity profiles form only when the difference in the granular temperature across the system exceeds a threshold value $`\mathrm{\Delta }T>\mathrm{\Delta }T_0`$. Second, in systems with large granular temperature differences $`\mathrm{\Delta }T>\mathrm{\Delta }T_0`$ at long times, the velocity profiles are linear near the ‘cold’ wall but highly nonlinear near the ‘hot’ wall. These profiles differ in shape from those found in the sheared but unvibrated systems xu . Finally, we point out that shear bands form when the average shear stress of the system $`\sigma _{xy}_t`$ falls below the yield stress $`\sigma _0`$ required to initiate sustained flow in a static system. Figs. 8 and 9 clearly show that differences in the mean-square velocity fluctuations across the system give rise to nonlinear velocity profiles. However, our results indicate that the granular temperature difference required to generate a nonlinear velocity profile must exceed a threshold value. For example, systems with vibration frequency $`\omega =1.8`$ (downward triangles) in Figs. 8 and 9 have relatively large granular temperature differences (slightly less than $`\mathrm{\Delta }T=10^3`$), but possess nearly linear velocity profiles. In contrast, systems with $`\mathrm{\Delta }T>10^3`$ (for example $`\omega =0.4`$, $`0.8`$, and $`1.4`$) possess strongly nonlinear velocity profiles. The threshold granular temperature difference, which is roughly equal to the granular temperature near the vibrated boundary, appears to agree with the glass transition temperature for a quiescent equilibrium system at the same average packing fraction remark . This correspondence seems reasonable since in equilibrium systems the temperature must exceed the glass transition temperature to cause large density fluctuations. Similarly, in sheared dissipative systems, it is difficult to create sufficiently large density gradients required for nonlinear velocity profiles if the granular temperature difference is below a threshold $`\mathrm{\Delta }T_0`$. We also find that the nonlinear velocity profiles that occur in the sheared and vibrated athermal systems have qualitatively different shapes compared to those found for the sheared but unvibrated systems. For example, the nonlinear velocity profiles for $`\omega =0.4`$ (squares) and $`1.4`$ (upward triangles) in Fig. 8 are composed of a linear portion that extends from the bottom boundary at $`y=0`$ to $`y0.8`$, and a strongly nonlinear part near the top boundary. We note that in vibrated systems the crossover between the linear and highly sheared behavior occurs where the local packing fraction switches from uniform to nonuniform. Finally, we will address an interesting claim made in Ref. varnik that shear bands form in sheared glassy systems when the average shear stress in the system falls below the yield shear stress required to induce flow in a static state. In systems that form shear bands, shear flow is confined to a small portion of the system while the remainder of the system remains nearly static. In contrast, all parts of the system flow when systems possess generic nonlinear velocity profiles. Does the constraint on the average shear stress guarantee that shear bands will occur in the repulsive athermal systems studied here? In Fig. 10, we compare the time-averaged shear stress $`\sigma _{xy}_t`$ sheared at fixed $`u`$ to the yield shear stress $`\sigma _0`$ required to induce flow in an originally static unsheared state at the same $`\varphi `$. We find that at small boundary velocities $`u`$, $`\sigma _{xy}_t`$ is slightly below $`\sigma _0`$. However, even though the average shear stress is below the yield stress, the velocity profiles are linear for long times (as shown in Fig. 3) and not highly localized or shear-banded. Thus, the condition $`\sigma _{xy}_t<\sigma _0`$ alone does not ensure that sheared repulsive athermal systems will form shear bands. An additional requirement must be satisfied to stabilize shear bands at long times—the systems must possess sufficiently large granular temperature differences. It should be noted, however, that the differences between the average and the yield shear stress are small in the systems we studied and that shear stress fluctuations may inhibit the formation of shear bands. Thus, more work should be performed to verify the presented results. In fact, we are now attempting to determine the variables that set the difference between $`\sigma _{xy}_t`$ and $`\sigma _0`$ and whether this difference persists in the large system limit. Fig. 11 shows that the average shear stress falls below the yield shear stress in the sheared and vibrated systems over a range of frequencies at $`u=0.00727`$ and $`0.0727`$, but not at $`0.727`$. We have also confirmed that the granular temperature differences in these systems satisfy $`\mathrm{\Delta }T>\mathrm{\Delta }T_0`$ over the range of frequencies $`0.4<\omega <1.4`$. Therefore, from the discussion above, one expects shear bands to form for the two higher shear rates, but not the lowest one. This prediction is confirmed in Fig. 12. At $`u=0.727`$, nonlinear velocity profiles form, but shear is not highly localized into a shear band. In contrast, at $`u=0.00727`$ there is a range of frequencies $`0.4<\omega <1.4`$ over which shear bands form. ## V Conclusions In this article we reported on recent simulations of model frictionless granular systems undergoing boundary-driven planar shear flow in 2D over a range of flow velocities and average densities, and for particles with varying degrees of inelasticity. These studies have produced several interesting and novel results that are relevant to a variety of jammed and glassy systems subjected to planar shear flow. First, we find that nonlinear velocity profiles are not stable at long times. Nonlinear velocity profiles form when the boundary velocity exceeds a characteristic speed set by the shear wave speed in the material, but they slowly evolve toward linear profiles at long times. In addition, the granular temperature and packing fraction profiles are initially spatially dependent but become uniform at long times. We measured the time $`t_l`$ required for the velocity profiles to become linear and for the granular temperature and density to become homogeneous throughout the system. We find that $`t_l`$ increases as $`u^{0.5}`$ at large $`u`$ and increases strongly as $`\varphi `$ approaches $`\varphi _{\mathrm{rcp}}`$ for nearly elastic systems. These results imply that sufficiently large and sustained granular temperature differences between the ‘hot’ and ‘cold’ boundaries are required to stabilize nonlinear velocity profiles at long times. We also studied systems in which vertical vibrations of the top or bottom boundary were superimposed onto planar shear flow to maintain granular temperature differences across the system. In the sheared and vibrated systems, we find that nonlinear velocity profiles are stable when the granular temperature difference exceeds a threshold value $`\mathrm{\Delta }T_0`$, which roughly corresponds to the glass transition temperature in an equilibrium system at the same average density. The nonlinear velocity profiles, however, differ in shape from those found previously for planar shear flow. The velocity profiles are composed of a linear part that exists where the packing fraction is spatially uniform and a nonlinear portion that exists where the packing fraction varies strongly in space. Finally, we have shown that the nonlinear velocity profiles become highly localized when the average shear stress in the system is below the yield shear stress. ## VI Future Directions Several additional studies are necessary to fully understand dense shear flows in granular systems, and these will be presented in future work xu2 . First, we intend to investigate the influence of dynamic and static friction forces on the long-time stability of nonlinear velocity profiles. Preliminary results suggest that dynamic friction is not sufficient to stabilize nonlinear velocity profiles in systems undergoing planar shear flow. In initial studies with weak dissipation and dynamic friction remark3 near $`\varphi _{\mathrm{rcp}}`$, we found that $`t_l`$ is similar to that for systems with dissipation and no dynamic friction. However, more extensive studies of dynamic as well as static friction are required to make a definitive statement about the influence of friction on the long-time stability of nonlinear velocity profiles aharonov ; volfson . Second, most experimental studies of velocity profiles in granular systems are performed in (angular) Couette, not planar shear cells. How does the geometry of the shear cell influence the velocity profiles? Are shear bands stable in frictionless, athermal systems undergoing Couette shear flow? Fig. 13 shows preliminary results from studies of 2D Couette shear flow in systems at $`\varphi =0.845`$, $`b_n=0.0375`$, and rotation rate $`\mathrm{\Omega }=0.01`$ in the counter clockwise direction. Panels (b), (c), and (d) show snapshots of the system $`2`$, $`6`$, and $`19`$ rotations after the initial configuration in panel (a). The snapshots provided in panels (a) - (c) reveal that the shear band in this system is approximately $`58`$ small particle diameters wide since the two highlighted particles closest to the outer boundary do not rotate significantly even after $`6`$ rotations of the inner boundary. A comparison of panels (c) and (d) shows that at long times the particles in the flowing region are able to diffuse perpendicular to the boundaries. The mean-square velocity fluctuations $`\delta v_\theta ^2`$ in the tangential direction and the tangential velocity $`v_\theta `$ normalized by the speed at the inner boundary are shown in Figs. 14 (a) and (b) as a function of the distance from the inner boundary $`(rR_1)/(R_2R_1)`$. Averages over varied numbers of rotations demonstrate that the nonlinear velocity profiles are stable at long times. As we found in our studies of planar shear flow, nonlinear velocity profiles occur when the granular temperature is nonuniform. However, in Couette shear flows, the shear stress is also spatially dependent. Thus, the distinct contributions from nonuniform shear stress and nonuniform granular temperature need to be disentangled. In future work, we will determine whether nonlinear velocity profiles persist when we vibrate the outer boundary to create a more uniform granular temperature profile. Finally, there have been several recent computational studies of effective temperatures defined from fluctuation dissipation relations, linear response theory, and elastic energy fluctuations in dense granular systems makse ; kondic ; ohernt . These studies have shown that a consistent effective temperature can be defined for dense shear flows, i.e. the effective temperatures $`T_{\mathrm{eff}}`$ obtained from the above definitions agree with each other but $`T_{\mathrm{eff}}`$ is much larger than the granular temperature of the system. Because $`T_{\mathrm{eff}}`$ describes fluctuations on long length and time scales, it is possible that the effective temperature will play a significant role in determining the shape of velocity profiles. Thus, it is important to study the effective temperature as well as the granular temperature in sheared glassy and athermal systems. #### Acknowledgments We thank J. Blawzdziewicz, A. Liu, and S. Nagel for helpful comments. Financial support from NSF DMR-0448838 (NX,CSO), NASA NNC04GA98G (LK), and the Kavli Institute for Theoretical Physics under NSF PHY99-07949 (CSO,LK) is gratefully acknowledged. We also thank Yale’s High Performance Computing Center for generous amounts of computing time.
warning/0506/math0506382.html	ar5iv	text	# Necessary And Sufficient Conditions For Existence of the LU Factorization of an Arbitrary Matrix. ## Notes about computation of LU factorization in floating point arithmetic . Every invertible matrix which has an LU factorization has a neighborhood in which every matrix is invertible and has LU factorization. Moreover, if we require that $`u_{ii}=1`$ then such a factorization is unique. Therefore we can define function $`A[L,U]`$. Such a function can be defined in some neighborhood of any invertible matrix that has an LU factorization. The function will be invertible and continuous in the neighborhood. The inverse will also be 1-1 and continuous. These facts make it possible to compute LU factorization of an invertible matrix satisfying conditions (1) using floating point arithmetic. However in general a matrix that has an LU factorization does not have a neighborhood in which every matrix has an LU factorization. Also if a factorization exists it does not have to be unique. Further LU factorization does not have to depend continuously on the entries of $`A`$. Thus, because of the possibility of the rounding error it is not generally possible to compute an LU factorization in the general case in floating point arithmetic using algorithm (1). ## Some Applications. ###### Theorem 6 Any n-by-n matrix $`A`$ can be written as $$A=U_1LU_2$$ with $`U_1`$ and $`U_2`$ upper triangular matrices and $`L`$ lower triangular matrix. Proof. By a series of type-3 elementary row operations matrix $`A`$ can be transformed into a matrix $`C`$ such that for all $`k=1,\mathrm{},n`$ we have $$rankC[\{1\mathrm{}k\}]=rankC[\{1\mathrm{}n\},\{1\mathrm{}k\}]$$ In particular, since $$rankC[\{1\mathrm{}k\},\{1\mathrm{}n\}]k$$ we see that $`C`$ satisfies conditions (1). In fact one can transform $`A`$ into $`C`$ using only elementary type-3 row operations that add multiples of $`i`$th row to the $`j`$th row with $`j<i`$. Any series of such operations can be realized by multiplying $`A`$ by an invertible upper triangular matrix $`U`$ on the left. We have $$C=UA$$ but by theorem (1) $$C=LU_2$$ with $`L`$ lower triangular and $`U_2`$ upper triangular. Since $`U`$ is invertible we can let $$U_1=U^1$$ That gives us $$U_1UA=A=U_1LU_2$$ $`\mathrm{}`$ ###### Note 8 Notice that $`U_1`$ can be taken invertible. Also because $$rankC[\{1\mathrm{}k\}]=rankC[\{1\mathrm{}n\},\{1\mathrm{}k\}]$$ it follows that $`L`$ can be taken invertible \[ see LM \]. ###### Corollary 2 Any n-by-n matrix $`A`$ can be written as $$A=L_1UL_2$$ with $`L_1`$ and $`L_2`$ lower triangular matrices and $`U`$ an upper triangular matrix. Proof. It follows from the fact that $`A^T`$ can be written as $$A^T=U_1LU_2$$ with $`U_1`$ and $`U_2`$ upper triangular matrices and $`L`$ lower triangular matrix. $`\mathrm{}`$ ###### Note 9 Notice that $`U`$ and $`L_2`$ can be taken invertible. We prove the following well known result ###### Theorem 7 Any n-by-n matrix $`A`$ can be written as $$A=PLU$$ with $`U`$ an upper triangular matrix, $`L`$ a lower triangular matrix and $`P`$ permutation matrix. Proof. Any matrix $`A`$ can be multiplied on the left by such permutation matrix $`P_0`$ that matrix $$C=P_0A$$ satisfies the following equality for all $`k=1,\mathrm{},n`$ $$rankC[\{1\mathrm{}k\}]=rankC[\{1\mathrm{}n\},\{1\mathrm{}k\}]$$ In particular, since $$rankC[\{1\mathrm{}k\},\{1\mathrm{}n\}]k$$ we see that $`C`$ satisfies conditions (1). therefore by theorem (1) $`C`$ can be written as $$C=LU$$ with a lower triangular $`L`$ and an upper triangular $`U`$. We let $$P=P_0^1$$ That gives us $$A=PP_0A=PC=PLU$$ $`\mathrm{}`$. ###### Corollary 3 Any n-by-n matrix $`A`$ can be written as $$A=LUP$$ with $`U`$ an upper triangular matrix, $`L`$ a lower triangular matrix and $`P`$ permutation matrix. Proof. It follows from the fact that $$A^T=P_0L_0U_0$$ with $`U_0`$ an upper triangular matrix, $`L_0`$ a lower triangular matrix and $`P_0`$ permutation matrix. Indeed, we can write $$A=LUP$$ with $$L=U_0^T$$ $$U=L_0^T$$ and $$P=P_0^T$$ $`\mathrm{}`$
warning/0506/cond-mat0506192.html	ar5iv	text	# The superheated Melting of Grain Boundary ## I Introduction The superheating has been found in a larger number of systems such as surface Carnevali ; Pluis ; Denier , small cluster SHVART , confined thin film ZHANGL1 and particles covered (or embedded in) by material with higher melting point Daeges . Generally, the melting of solid material is heterogeneous process with the nucleation mechanism at surfaces or interfaces Cahn ; Madd . Providing heterogeneous nucleation could be avoided by means of suitable coating Daeges or internal heating Khaikin , the metal crystal can be in the superheated state, its melting is completed by a thermodynamically instability resulting in homogeneous disordering and catastrophic mechanism with the stability limit from 0.2$`T_m`$ to 2.0$`T_m`$ Lele ; Kauzmann ; Fecht ; Tallon ; Lu . A large number of researches have been contributed to the role of surface for the melting of crystal. The superheated melting of fcc(110) and fcc(100) surfaces are virtually never observed Frenken ; Tolla ; Sun . The only example of the superheated surface is the small crystal strictly confined by high-symmetry fcc(111) facts Carnevali ; Pluis ; Denier . GB as another important quasi-2D defect also leads to the heterogeneous melting of solid material. Quite a number of studies have shown that GB can’t melt below the $`T_m`$ (not premelting) Nguyen ; Kikuchi ; Ciccotti ; Nguyen1 ; Lutsko ; Lutsko1 ; Broughton ; Plimpton ; Carrion ; Plimpton1 ; Fan ; Phillpot . Using the MD simulation, Kikuchi and Cahn, Ciccotti $`et.al.`$ showed that GB doesn’t melt until temperature reaches to melting point of bulk. Nguyen $`et.al.`$ using the more accurate interatomic potential by embedded-atom methods(EAM) studied the high-temperature GB structure, they found that, close to melting point T<sub>m</sub>, the GB structure was disordered, quite liquid-like and meta-stable, and over a long interval of simulation the underlying crystalline order can re-emerge. The experiment (T.E.Hsieh and R. B. Balluffi) using the HREM(High Resolution Electron Microscopy) methods supported above arguments and showed that aluminum GB did not melt below 0.999$`T_m`$ Hsieh . It is more surprising that some the simulations also hint that some high symmetric GBs similar to high symmetric surface can probably melt by the superheated Broughton . In this work, parallel to Di Tolla’s work Tolla for surface we study the possibility of superheated high-symmetry GB by a theoretical model and MD simulation of a symmetric aluminum GB. When temperature beyond melting point a crystal reaches the superheated state, all liquid nuclei must be smaller than a critical size. These liquid nuclei are unstable and able to re-crystallize again. The liquid nuclei easily form at surfaces, grain boundary and other solid defects regions. So, to avoid larger liquid nucleus, the crystal must be prepared with the lowest number of solid defects. It’s easier to study the superheating of crystal in computer simulation than in experiment. We can construct prefect crystal using the periodic boundary condition in computer simulation. In experiment it’s very difficult to obtain infinite volume prefect crystal. We can’t eliminate the influence of surfaces, grain boundaries, dislocations and other defects with complicated structures. All these defects have potentially become the liquid nuclei to melt crystal. However surface effects can partially remove by coating other material with higher melting point or internally heating the material. By these methods we can obtain superheated crystalline grains. Additionally, the grain boundary itself may probably become the liquid nucleus near the melting point of crystal. Based on our theoretical model, it’s possible that, under condition of the absent of critical nucleating cores, grain boundary doesn’t melt even the temperature beyond melting point. We hope to find the superheated grain boundary in computer simulation, although it is difficult to find in experiments due to different kind of unavoidable nucleation mechanism. Our simulation will show that high-symmetry grain boundary can sustain above the melting point of crystal without other nucleation mechanism. The superheated grain boundary is easily understood by proximity effects if we consider the grain boundary is sandwiched between two superheated crystalline grains, while the superheated grains can be obtained by properly internal heating. ## II The Model of GB Melting The melting point $`T_m`$ of a solid may be defined as the temperature with the coexistence of solid phase and liquid phase. For a solid with surfaces or grain boundaries, the melting is generally completed by the mechanism of heterogeneous nucleation. We consider that a liquid film with thickness $`2l`$ forms between two semi-infinite solid (Fig. 1(a)). The change of free energy per unit area is taken as $$\mathrm{\Delta }F(l)=2\rho Ll(1T/T_m)+\mathrm{\Delta }\gamma (l)$$ (1) where $`\rho `$ is the liquid density, $`L`$ the latent heat of melting, $`\mathrm{\Delta }\gamma (l)`$ the difference between the overall free energy $`\gamma _{SLSL}`$ of two interacting solid-liquid interfaces separated by a distance $`l`$ and the GB energy per unit area $$\mathrm{\Delta }\gamma (l)=\gamma _{SLSL}\gamma _{GB}$$ (2) By extending the Cahn’s wetting theory to solid-solid interface Gennes , using $`\mathrm{\Delta }\gamma (0)`$ = 0 and only considering the short range interaction, we may obtain $$\mathrm{\Delta }\gamma (l)=\mathrm{\Delta }\gamma _{\mathrm{}}(1e^{l/\xi }).$$ (3) where $`\mathrm{\Delta }\gamma _{\mathrm{}}=2\gamma _{SL}\gamma _{GB}`$ is the difference of the interface energy of two isolated solid-liquid interface $`\gamma _{SL}`$ and the GB energy $`\gamma _{GB}`$, $`\xi `$ is the width of solid-liquid interface. The condition of GB melting is defined by following process. We may image there is a droplet in GB region (Fig. 1(b)), the equilibrium condition of the droplet is $$\gamma _{GB}2\gamma _{SL}cos(\theta )=0$$ (4) where $`\theta `$ is the wetting angle. If $`\mathrm{\Delta }\gamma _{\mathrm{}}>0`$, $`i.e.`$ $`\gamma _{GB}<2\gamma _{SL}`$, $`\theta `$ is a finite value, the droplet can survive in the region of GB (partial wetting), the whole GB can’t be wetted. As $`\mathrm{\Delta }\gamma _{\mathrm{}}<0`$, $`i.e.`$ ,$`\gamma _{GB}>2\gamma _{SL}`$, $`\theta `$ can’t be defined. A liquid film forms in GB region (Wetting) and the GB melts. Thus $`\mathrm{\Delta }\gamma _{\mathrm{}}=0`$, $`i.e.`$ $`\gamma _{GB}=2\gamma _{SL}`$ and $`\theta =0`$, may be considered as the criteria of GB Melting. In early time, one was decline to think that GBs melt below melting temperature, that is, at certain temperature below $`T_m`$, $`\mathrm{\Delta }\gamma _{\mathrm{}}<0`$. In this paper we only consider the possibility of GB superheated melting, that is, as $`T>T_m`$, $`\mathrm{\Delta }\gamma _{\mathrm{}}>0`$. For superheated melting, the insert of Fig. 1(c) show that $`\mathrm{\Delta }F(l)`$ has a local minimum at $`l=0`$, and approaches to negative infinite as $`l`$ approach infinite. There is a maximum at $`l_c`$, which is $$l_c=\xi Ln(\frac{\mathrm{\Delta }\gamma _{\mathrm{}}T_m}{2L\rho \xi (TT_m)})$$ (5) The Fig. 1(c) shows the temperature dependence of $`l_c`$. According to Fig. 1(c), at a certain temperature $`T(>T_m`$) and $`l<l_c`$, the system can reduce the its free energy by decreasing the thickness of liquid film until the thickness reaches to zero, that is, the system crystallizes. When $`l>l_c`$, the system can reduce its free energy by increasing the thickness of liquid film until the thickness reach to infinite, the GB melts. $`l_c`$ is the critical thickness at $`T`$. If at a temperature $`l_c=0`$, any small thickness can lead to the melting of GB. The temperature is named as the maximum superheated temperature $`T_s`$, expressed as $$T_s=T_m(1+\frac{\mathrm{\Delta }\gamma _{\mathrm{}}}{2L\rho \xi })$$ (6) $`l_c`$ is also be expressed as $$l_c=\xi Ln(\frac{T_sT_m}{TT_m})$$ (7) The critical nucleus is extremely large for $`TT_m`$ and reduces rapidly with the increase of the superheating degree. At $`T_s`$, $`l_c=0`$ (Fig. 1(c)) and the spontaneous melting happens. Between $`T_m`$ and $`T_s`$ the superheated states is meta-stable, although the melting doesn’t occur. In following several sections, by MD simulations of Aluminum GB melting and above model, we prove that GB can preserve crystalline even above $`T_m`$ until temperature reaches to the maximum superheated temperature $`T_s`$. Our simulation shows the behavior of superheating GB, that is, for $`T_m<`$ $`T`$ $`<T_s`$, there exists a critical width $`l_c`$ of liquid film. When the width of the artificially added liquid is larger than $`l_c`$ the GB melts, or the liquid film will reduce and the effect of crystallization is dominating. ## III The Molecular Simulation of Gain Boundary Molecular Dynamic simulations of crystal JIN and high-symmetry surface Tolla have shown the superheating crystal without critical nuclei such as point defects and liquid drops. By internal heating, crystal with high-symmetry fcc(111) surface will be superheating. This is because the process of nucleation is homogeneous in the surface region. For some low symmetry surfaces such as fcc(110) surface, the anisotropy leads to heterogeneous nucleation and high concentration of defects, superheating is extremely difficult to achieve. Reconstruction and roughening are two main structural transitions for a surface when increasing temperature. Behavior of grain boundary is more complicated than surface, including the migration, bending, sliding, zigzag and faceting transition. Our simulations show that it is difficult to control the homogeneous nucleation in computer simulation to obtain superheated grain boundary. However we can use symmetric grain boundary to make atoms homogenously distribute in GB region. Instead of periodic boundary, we have fixed two boundaries of simulation cell parallel to GB plane to decrease the possibility of grain-boundary sliding. In this work we choose Al symmetric $`\mathrm{\Sigma }`$13 (320) tilt grain boundary Fig 1(d) as our simulation cell with mis-orientation $`67.8^{}`$, the tilt axis of the boundary is along the Z direction, the fixed boundary-conditions is used in the X direction perpendicular to GB plane, and periodic boundary condition are used in the Y and Z directions. The widths along X, Y, Z are 100$`\AA `$, 43$`\AA `$, 12$`\AA `$ respectively. Interatomic potential plays the most important role in molecular dynamics simulations. In this paper, by using a more realistic potential, a Glue potential developed by F. Ercolessi and J. B. Adams E1 , The glue potential has been used in a large number of the simulations of surface, cluster, liquid, and crystal, and the simulation results are perfectly consistent with experimental results Carnevali ; Tolla ; Shu ; E3 ; E4 . The lattice constant $`a_0`$ for Al at 0K is 4.032$`\AA `$. For the glue potential the melting points is about 936K E2 , which is close to the experimental melting point of Aluminum (about 933K). The MD simulations are carried out at constant temperature, constant volume and constant atomic number. The time step is 0.06 ( about 0.003psec ), at each temperature the runs are made about 20000 time steps ( about 60psec ). The static structure-factor $`S(𝐊)`$ for specific reciprocal-space vector $`𝐊`$ (Eq. 8) represents a quantitative measure for the long-range order and can be used as order parameter describing the transition between disorder and order. $$S(𝐊)=<\underset{iGB}{}\mathrm{exp}(i𝐊𝐫_i(t))^2>/N_{GB}^2$$ (8) $`𝐫_i(t)`$ is the position of $`i`$th atom at time $`t`$. $`<\mathrm{}>`$ is indicative of the time average. $`N_{GB}`$ is the number of atoms in grain-boundary region. For a crystal the $`S(𝐊)`$ is approaching 1, for liquid it approaching to 0. In order to study the stability and the nature of disordered grain-boundary, we define a distribution function of $`P(S)`$ which is the statistics of the value of $`S(𝐊)`$ ($`𝐊=\frac{2\pi }{a_0}(0,0,1))`$ of all molecular-dynamics time steps. Fig. 2 shows the distribution functions at several temperatures from 400K to 1050K. The centers of the peaks represent the degree of disorder and order, the widths of the peaks as a measure of the fluctuation of $`S(𝐊)`$ represent the instability of the GB structure. From Fig. 2, the peak is very narrow and the center of the peak is very near 1 in low temperature regime ($`T`$ 400K) , and this implies that the structure of grain boundary is very order and stable in low temperature regime; At 850K, the peak is board ( the minimum is at about 0.5 and the maximum is at about 0.9 ) and the center of the peak is much less than 1 , which implies that the GB structure is disordered in this temperature. At about 950K, the width of peak is extremely board (the minimum reaches 0.1 and the maximum still retains at about 0.9), the fact implies the structure of grain boundary is rather unstable, sometime the structure of grain boundary is rather disorder like liquid because of the $`S(𝐊)`$ approaching 0, and sometimes it just likes a crystal structure with long range order because of the $`S(𝐊)`$ probably approaching 1. This shows that at this temperature the coexistence of liquid phase and solid phase is reached. Above results show that the melting-point of this system is very close to 950K. In our simulation, the large $`S(𝐊)`$ fluctuation near 950K just indicates the signal of solid-liquid phase transition. At 1050K the peak moves to the left and become narrow. The position of the peak is close to zero, the GB is melting. Our results also show that the maximum superheated temperature $`T_s`$ is close to 1050K. We also calculated the pair correlation functions $`g(r)=\frac{1}{4\pi \rho N}<_{ij}\delta (rr_{ij})>`$ at various temperatures which show the liquid behavior at 1050K. Fig. 3(a,b) illustrate the GB structures at 850K and 1050K. Our model shows that between $`T_m`$ and $`T_s`$, GB will enter a new superheated state of GB. The new state is characterized by (1) the coexistence of liquid and solid; (2) the smaller size of liquid nucleus than that of the critical nucleus at that temperature prevents from the melting of GB although $`T>T_m`$. Fig. 3(c,d) show the competition of liquid and solid phase at 975K in superheated state. Sometimes there exists a liquid-like layer in the grain-boundary region (c) but it is meta-stable and may disappear and crystalline phase re-emerges at following time steps (d). We will show that the superheated state (975K) is rather different from the high-temperature disordered state at 850K. A liquid layer as the melting nucleus can be artificially added to the GB region by following methods: At a certain temperature $`T`$, we sample some layers with width $`2l`$, the atoms in these layers are heated up to an appropriate temperature $`T_a`$($`>T`$) until a liquid layer forms. By allowing all atoms relaxation at temperature $`T`$ again, we can obtain a new equilibrium structure at temperature $`T`$. Fig. 4 shows that both the initial GB configurations Fig. 4(a,c,e) having already added a liquid film with widths $`2l`$ and the final equilibrium configurations Fig. 4(b,d,f) at 975K and 850K respectively. At 850K, for $`2l=30\AA `$, the liquid layers disappears after the relaxation about 50ps Fig. 4(a,b). However at 975K for $`2l=20\AA `$, the liquid layer disappears Fig. 4(c,d)and for $`2l=30\AA `$ the layer of liquid is still existent Fig. 4(e,f) after relaxation about 50ps. Therefore, we can obtain $`20\AA `$ $`<2l_c<`$ $`30\AA `$ at $`T`$=975K. Above results also show that the superheated state (975K) is very different in nature from the high temperature disorder state (850K). For high temperature disorder state, the liquid film can’t induce the melting of GB, but for superheated GB state, the GB melts only when the width of liquid film is larger than a critical width $`l_c`$. In order to define the correlation length $`\xi `$ and the thickness of liquid film, we calculate the atom density profile corresponding Fig. 4(f). Our results show that at 975K, the critical width $`2l_c`$ is between $`20\AA `$ and $`30\AA `$ and $`\xi =10\AA `$, thus $`1.0<l_c/\xi <1.5`$, which is consistent with the theory model $`l_c/\xi =1.4`$ with $`T_m=950K`$ and $`T_s=1050K`$ in Eq. 7. ## IV Discussion and Conclusion In summary, we study the possibility of superheated GB state by both a theoretical model and MD simulation. Our results indicate that we can obtain superheated grain boundary by having properly controlled homogeneously nucleating precession when increasing temperature. If there are liquid nuclei whose sizes are larger than a critical size the superheated grain boundary melts. Or the grain boundary waits for homogeneously melting when temperature higher than maximum superheated temperature. We must justify that pressure plays important roles for the superheated in experiments and computer simulations. In experiments, both coating with high-melting-point materials and heating internally induce the internal pressure in melting region. In our simulation also there is internal pressure in melting region because the size of simulation cell doesn’t change companying with the increasing temperature. In this paper we only consider the superheating state due to the nucleation mechanism. The pressure mechanism and nucleation mechanism become intertwined and influence the melting of materials. Pressure lead to the increase of melting points. The melting point is about 950K in our simulation and higher than the experimental melting points T<sub>exp</sub>=933K and T<sub>glue</sub>=936K E2 in simulation using glue potential without internal pressure. Pressure increases melting point less than T<sub>m</sub>-T$`{}_{glue}{}^{}`$950K-936K=14K. However the nucleation mechanism leads to the superheating about T<sub>s</sub>-T$`{}_{m}{}^{}`$100K. The proximity effects of the superheated grains are also important to induce the superheated grain boundary sandwiched between two properly superheating grains. Because superheated state is meta-stable state in phase diagram, we don’t expect it’s long-live. The superheated materials melt by the nucleation mechanism or homogeneously melt at higher temperature. ## Acknowledgements Author is greatly indebted to professor D.Y.Sun and Yizhen He for valuable suggestions. This work is financially supported by Laboratory of Internal Friction and Defects in Solids, Chinese Academy of Sciences (currently named as Key Laboratory of Materials Physics, Chinese Academy of Sciences), the National Nature Science Foundation of China and Chinese Academy of Sciences under KJCX2-SW-W11.
warning/0506/gr-qc0506042.html	ar5iv	text	# On the theory of the skewon field: From electrodynamics to gravity ## 1 Introduction At the centennial of the proposition of special relativity theory by Einstein (1905), it is worthwhile to remember that Einstein’s paper was “On the electrodynamics of moving bodies,” see . The task Einstein had taken up was to develop a consistent framework for accommodating Maxwell’s theory of electrodynamics as well as classical mechanics. The tool for achieving this was to study the motion of charged bodies under the action of an electromagnetic field. Maxwell’s theory, suitably interpreted, survived, Newton’s mechanics had to be extended if high relative velocities were involved. Thus, Maxwell’s theory emerged as a prime example of a special relativistic field theory that is intrinsically related to the Poincaré group (also known as inhomogeneous Lorentz group). Accordingly, electrodynamics is conventionally thought to take place in the flat Minkowski space of special relativity as pointed out by Minkowski in his geometrical formulation of special relativity in 1908, see . The success of special relativity was so striking that the historical fact of the close association of electrodynamics with special relativity stuck in the minds of most physicists and is believed to be a physical fact — even though the development of classical electrodynamics during the last 100 years shows the opposite: The foundations of electrodynamics have nothing to do with special relativity and the Poincaré group, they are rather of a generally covariant (“topological”) nature based on the conservation laws of electric charge and magnetic flux. This development started with Einstein who shortly after the publication of his general relativity theory observed that the Maxwell equations can be formulated in such a way that neither the metric nor the Christoffel symbols appear in them. In his notation ($`\mu ,\nu ,\mathrm{}=0,1,2,3`$), they read<sup>5</sup><sup>5</sup>5Einstein used subscripts for denoting the coordinates $`x`$, i.e., $`x_\tau `$ etc. Moreover, we dropped twice the summation symbols $`\mathrm{\Sigma }`$. $$\frac{F_{\rho \sigma }}{x^\tau }+\frac{F_{\sigma \tau }}{x^\rho }+\frac{F_{\tau \rho }}{x^\sigma }=0,^{\mu \nu }=\sqrt{g}g^{\mu \alpha }g^{\nu \beta }F_{\alpha \beta },\frac{^{\mu \nu }}{x^\nu }=𝒥^\mu .$$ (1) Here we draw on our paper . The Maxwell equations (1)<sub>1</sub> and (1)<sub>3</sub> are apparently metric free. Moreover, since the excitation $`^{\mu \nu }`$ is considered to be a tensor density of type $`\left[\genfrac{}{}{0pt}{}{2}{0}\right]`$ and the field strength $`F_{\rho \sigma }`$ a tensor of type $`\left[\genfrac{}{}{0pt}{}{0}{2}\right]`$, both equations are — even though only partial derivatives operate in them — covariant under general coordinate transformations (diffeomorphisms). In other words, the system consisting of (1)<sub>1</sub> and (1)<sub>3</sub> doesn’t couple to the gravitional potential as long as (1)<sub>2</sub> is not substituted. A similar presentation of Maxwell’s equations was given by Einstein in his “Meaning of Relativity” in the part on general relativity. At first sight, this separation of the Maxwell equations into the Maxwell equations proper of (1)<sub>1</sub> and (1)<sub>3</sub> and the spacetime relation (1)<sub>2</sub> may appear to look like a formal trick. However, it is well known that the electric excitation $`𝒟`$ and the magnetic excitation $``$ are both directly measurable quantities, see, e.g., Raith . Accordingly, the electromagnetic field is represented operationally not only by means of the electric and magnetic field strengths $`E`$ and $`B`$ — both measured via the Lorentz force — but also by means of the electric and magnetic excitations $`𝒟`$ and $``$. On the foundations of general relativity, besides the equivalence principle, there lays the principle of general covariance. And the Maxwell equations (1)<sub>1</sub> and (1)<sub>3</sub> are generally covariant and metric independent. Since in general relativity the metric $`g`$ is recognized as gravitational potential, it is quite fitting that the fundamental field equations of electromagnetism do not contain the gravitational potential. Consequently, the Maxwell equations in their premetric form are valid in any 4D differential manifold, provided the latter can be split locally into 1+3. Accordingly, they are not only beyond special relativity, but also beyond general relativity. The point of view that the fundamental structure of electrodynamics can be best understood because of the existence of conservation laws that can be formulated generally covariant and metric-free has been mainly developed by Kottler (1922), É.Cartan (1923), and van Dantzig (1934), see . Modern presentations of this “premetric electrodynamics” have been given, e.g., by Truesdell-Toupin , Post , Kovetz , Rubilar , Hehl & Obukhov , Kiehn , Delphenich , and Lindell<sup>6</sup><sup>6</sup>6Lindell’s presentation of electrodynamics in the framework of exterior differential forms is metric independent. However, in order to make himself understood to his engineering public, he often interpretes the differential-form expressions in terms of metric-dependent Gibbsian vector expressions. , see also Itin . ## 2 Premetric electrodynamics ### 2.1 The Maxwell equations The conservation of electric charge leads to the inhomogeneous Maxwell equation: $$dIH=J(_j\stackrel{ˇ}{}^{ij}=\stackrel{ˇ}{J}^i).$$ (2) The first version in exterior calculus is written in terms of the twisted excitation 2-form $`IH=IH_{ij}dx^idx^j/2`$ and the twisted current 3-form $`J=J_{ijk}dx^idx^jdx^k/6`$. The translation into the corresponding version in components is achieved by $`\stackrel{ˇ}{H}{}_{}{}^{ij}:=ϵ^{ijkl}IH_{kl}/2`$ and $`\stackrel{ˇ}{J}{}_{}{}^{i}:=ϵ^{ijkl}J_{jkl}/6`$, where $`ϵ^{ijkl}`$ is the totally antisymmetric Levi-Civita symbol with components of value $`\pm 1,0`$. Magnetic flux conservation is represented by the homogeneous Maxwell equation Fig.1. The tetrahedron of the electromagnetic field. The excitation $`IH=(,𝒟)`$ and the field strength $`F=(E,B)`$ are 4-dimensional quantities of spacetime, namely 2-forms with and without twist, respectively. They describe the electromagnetic field completely. Of electric nature are $`𝒟`$ and $`E`$, of magnetic nature $``$ and $`B`$. In 3 dimensions, $``$ and $`E`$ are twisted and untwisted 1-forms, respectively; analogously, $`𝒟`$ and $`B`$ are twisted and untwisted 2-forms, respectively. The magnetic and the electric excitations $`IH=(,𝒟)`$ are extensities, also called quantities (how much?), the electric and the magnetic field strengths $`E`$ and $`B`$ are intensities, also called forces (how strong?). $$dF=0(_{[i}F_{jk]}=0),$$ (3) with the field strength 2-form $`F=F_{ij}dx^idx^j/2`$. The decompositions into time and space read $$IH=dt+𝒟,F=Edt+B,$$ (4) compare the scheme in Fig.1. Conservation laws can be reduced to counting procedures. No distance concept is required in this context, rather only the ability to circumscribe a definite volume or an area. As a consequence, no metric occurs anywhere in (2), (3), and (4). ### 2.2 Local and linear spacetime relation Excitation and field strength in vacuum (generalization to media is possible) are assumed to be related by a local and linear relation $$IH=\kappa [F](IH_{\alpha \beta }=\frac{1}{2}\kappa _{\alpha \beta }{}_{}{}^{\gamma \delta }F_{\gamma \delta }^{}).$$ (5) The constitutive tensor $`\kappa _{\alpha \beta }{}_{}{}^{\gamma \delta }=\kappa _{\beta \alpha }{}_{}{}^{\gamma \delta }=\kappa _{\alpha \beta }^{\delta \gamma }`$ has 36 independent components. These components can be understood in terms of the tensor-valued 2-form $$𝔎^{\alpha \beta }=\frac{1}{2}\kappa _{\gamma \delta }{}_{}{}^{\alpha \beta }\vartheta _{}^{\gamma }\vartheta ^\delta .$$ (6) Then (5)<sub>1</sub> can be written explicitly as $$IH=\frac{1}{2}𝔎^{\alpha \beta }e_\beta e_\alpha F.$$ (7) The constitutive 2-form $`𝔎^{\alpha \beta }`$ can be split, according to 36=20+15+1, into three irreducible pieces, $$𝔎^{\alpha \beta }=\underset{\text{principal}}{\underset{}{{}_{}{}^{(1)}𝔎_{}^{\alpha \beta }}}+\underset{\text{skewon}}{\underset{}{{}_{}{}^{(2)}𝔎_{}^{\alpha \beta }}}+\underset{\text{axion}}{\underset{}{{}_{}{}^{(3)}𝔎_{}^{\alpha \beta }}}(\kappa _{\alpha \beta }{}_{}{}^{\gamma \delta }=\underset{A=1}{\overset{3}{}}{}_{}{}^{(A)}\kappa _{\alpha \beta }^{}{}_{}{}^{\gamma \delta }).$$ (8) A detailed proof is given in the Appendix. In particular, $${}_{}{}^{(2)}𝔎_{}^{\alpha \beta }=\overline{)}𝔎^{[\alpha }\vartheta ^{\beta ]},{}_{}{}^{(3)}𝔎_{}^{\alpha \beta }=\alpha \vartheta ^\alpha \vartheta ^\beta .$$ (9) If the spacetime relation can be derived from a Lagrangian, then the skewon piece $`{}_{}{}^{(2)}𝔎_{}^{\alpha \beta }`$ has to vanish. On the other hand, $`{}_{}{}^{(2)}𝔎_{}^{\alpha \beta }`$ is a permissible structure provided it is related to dissipative processes. This is indeed the case, see . The hypothesis of the existence of $`{}_{}{}^{(2)}𝔎_{}^{\alpha \beta }`$ was proposed by three of us . Its effect on the light propagation has been studied in the meantime . The axion piece $`{}_{}{}^{(3)}𝔎_{}^{\alpha \beta }`$ had been proposed much earlier in an elementary particle context, see the axion electrodynamics of Wilczek and the literature given there. In particular for a comparison with the literature it is convenient to introduce the equivalent constitutive tensor density $$\chi ^{\alpha \beta \gamma \delta }:=\frac{1}{2}ϵ^{\alpha \beta \mu \nu }\kappa _{\mu \nu }{}_{}{}^{\gamma \delta },\chi ^{\alpha \beta \gamma \delta }=\underset{A=1}{\overset{3}{}}{}_{}{}^{(A)}\chi _{}^{\alpha \beta \gamma \delta },$$ (10) with 36=20+15+1 independent components. Its skewon and its axion pieces can be mapped to a tensor (15 components) and a pseudo-scalar (1 component), respectively, $$\overline{)}S_\alpha {}_{}{}^{\beta }:=\frac{1}{4}\widehat{ϵ}_{\alpha \gamma \delta ϵ}^{(2)}\chi ^{\gamma \delta ϵ\beta },\alpha :=\frac{1}{4!}\widehat{ϵ}_{\alpha \beta \gamma \delta }^{(3)}\chi ^{\alpha \beta \gamma \delta }.$$ (11) We have $`\overline{)}S_\alpha {}_{}{}^{\alpha }=0`$. For the 1-form $`\overline{)}S^\alpha :=\overline{)}S_\beta {}_{}{}^{\alpha }\vartheta _{}^{\beta }`$, with $`e_\alpha \overline{)}S^\alpha =0`$, we find by some algebra, $$\overline{)}S^\alpha =\frac{1}{2}\overline{)}𝔎^\alpha ,\alpha =\frac{1}{12}𝔎.$$ (12) Up to here, our considerations were premetric. If we put the physical dimension of $`{}_{}{}^{(1)}𝔎_{}^{\alpha \beta }`$ into a function $`\lambda (x)`$, the so-called dilaton, then we have 1 dilaton field $`\lambda `$, 19 remaining components of the principal part $`{}_{}{}^{(1)}𝔎_{}^{\alpha \beta }/\lambda `$, 15 skewon components $`\overline{)}S_\alpha ^\beta `$, and 1 axion component $`\alpha `$. This is as far as we can go with the premetric concept. The study of the properties of light propagation can help to constrain the constitutive tensor of spacetime. The skewon field, a specific kind of permeability/permittivity of spacetime, will be in the center of our interest. It is non-Lagrangian and dissipative, and it diffracts light. Here we want to address the problem of how it could couple to gravity. However, first we want to turn our attention to its electromagnetic energy-momentum. ## 3 The electromagnetic energy-momentum density of the skewon field The 3-form of the electromagnetic energy-momentum current can be taken from , e.g.: $$\mathrm{\Sigma }_\alpha :=\frac{1}{2}[F(e_\alpha IH)IH(e_\alpha F)].$$ (13) We decompose the 3-form according to $`\mathrm{\Sigma }_\alpha =\mathrm{\Sigma }_{klm\alpha }dx^kdx^ldx^m/6`$. Then the corresponding energy-momentum tensor in tensor calculus can be defined as $`𝒯_i{}_{}{}^{j}:=ϵ^{jklm}\mathrm{\Sigma }_{klmi}/6`$ or $$𝒯_i{}_{}{}^{j}=\frac{1}{4}\delta _i^jF_{kl}\stackrel{ˇ}{}^{kl}F_{ik}\stackrel{ˇ}{}^{jk}.$$ (14) Let us now turn to the skewon part. The excitation can be decomposed in principal, skewon, and axial parts according to $`IH={}_{}{}^{(1)}IH+{}_{}{}^{(2)}IH+{}_{}{}^{(3)}IH`$. Since $`IH`$ enters (13) linearly, we find $`\mathrm{\Sigma }_\alpha ={}_{}{}^{(1)}\mathrm{\Sigma }_{\alpha }^{}+{}_{}{}^{(2)}\mathrm{\Sigma }_{\alpha }^{}+{}_{}{}^{(3)}\mathrm{\Sigma }_{\alpha }^{}`$, with $`\mathrm{\Sigma }_\alpha \|_{\mathrm{skewon}}{}_{}{}^{(2)}\mathrm{\Sigma }_{\alpha }^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[F(e_\alpha {}_{}{}^{(2)}IH){}_{}{}^{(2)}IH(e_\alpha F)]`$ (15) $`=`$ $`{\displaystyle \frac{1}{2}}e_\alpha (F{}_{}{}^{(2)}IH){}_{}{}^{(2)}IHe_\alpha F.`$ The skewon part of the excitation was derived in (75) as $${}_{}{}^{(2)}IH=\frac{1}{2}\overline{)}𝔎^\alpha e_\alpha F.$$ (16) We multiply it by $`F`$, apply the anti-Leibniz rule for the interior product, and recall that $`e_\alpha \overline{)}𝔎^\alpha =0`$: $`F{}_{}{}^{(2)}IH`$ $`=`$ $`{\displaystyle \frac{1}{2}}F\overline{)}𝔎^\alpha e_\alpha F={\displaystyle \frac{1}{2}}e_\alpha (F\overline{)}𝔎^\alpha F)+{\displaystyle \frac{1}{2}}(e_\alpha F)\overline{)}𝔎^\alpha F`$ (17) $`+{\displaystyle \frac{1}{2}}F(e_\alpha \overline{)}𝔎^\alpha )F={\displaystyle \frac{1}{2}}F\overline{)}𝔎^\alpha e_\alpha F.`$ Thus, $`F{}_{}{}^{(2)}IH=0`$. We substitute this into (15) and find $$\mathrm{\Sigma }_\alpha \|_{\mathrm{skewon}}=(e_\alpha F)\overline{)}S^\beta (e_\beta F)\text{(premetric result)}.$$ (18) A similar computation can be performed for the energy-momentum tensor. We have (, p.256) $$^{(2)}IH_{ij}=2\overline{)}S_{[i}{}_{}{}^{k}F_{j]k}^{}\text{or}{}_{}{}^{(2)}\stackrel{ˇ}{}_{}^{mn}=ϵ^{mnij}\overline{)}S_i{}_{}{}^{k}F_{jk}^{}.$$ (19) On substitution of this into (14), we find the skewon part of the energy-momentum tensor as<sup>7</sup><sup>7</sup>7If we translate (20) into the energy-momentum 3-form, then, in components, the corresponding formula reads: $$\mathrm{\Sigma }_{ijk\alpha }\|_{\mathrm{skewon}}=\frac{3}{4}\left(\overline{)}\kappa _\alpha {}_{}{}^{m}F_{m[i}^{}F_{jk]}+\overline{)}\kappa _{[i}{}_{}{}^{m}F_{jk]}^{}F_{\alpha m}+2\overline{)}\kappa _{[i}{}_{}{}^{m}F_{j\|m}^{}F_{\alpha \|k]}\right).$$ . $$𝒯_i{}_{}{}^{j}\|_{\mathrm{skewon}}^{}=ϵ^{jklm}F_{ik}\overline{)}S_l{}_{}{}^{n}F_{nm}^{},$$ (20) which is clearly equivalent to (18) and also of premetric nature. ### 3.1 Trace We transvect (15) with the coframe and recall that $`F{}_{}{}^{(2)}IH=0`$: $$\vartheta ^\alpha {}_{}{}^{(2)}\mathrm{\Sigma }_{\alpha }^{}=\vartheta ^\alpha {}_{}{}^{(2)}IHe_\alpha F=2F^{(2)}IH=0.$$ (21) Thus, the tracelessness is proved: $$\vartheta ^\alpha {}_{}{}^{(2)}\mathrm{\Sigma }_{\alpha }^{}=0\text{(premetric result)}.$$ (22) Equivalently, $`𝒯_i{}_{}{}^{i}=0`$. Thus, the skewonic part of the energy-momentum is tracefree. Since $`\vartheta ^\alpha \mathrm{\Sigma }_\alpha =0`$, we find the analogous property for the principal part: $`\vartheta ^\alpha {}_{}{}^{(1)}\mathrm{\Sigma }_{\alpha }^{}=0`$. ### 3.2 Antisymmetric part For these considerations we need the existence of a metric. We lower the index of the coframe $`\vartheta _\alpha :=g_{\alpha \beta }\vartheta ^\beta `$, multiply the energy-momentum from the left, and antisymmetrize: $$\vartheta _{[\alpha }{}_{}{}^{(2)}\mathrm{\Sigma }_{\beta ]}^{}={}_{}{}^{(2)}IH\vartheta _{[\alpha }e_{\beta ]}F=\vartheta _{[\alpha }(e_{\beta ]}F)\overline{)}S^\gamma (e_\gamma F).$$ (23) This obviously does not vanish. In contrast, in conventional Maxwell-Lorentz vacuum electrodynamics, we have, of course, $`\vartheta _{[\alpha }\mathrm{\Sigma }_{\beta ]}=0`$, that is, a symmetric energy-momentum, see . Alternatively, one can consider the 2-form $`W:=e^\alpha ^{(2)}\mathrm{\Sigma }_\alpha `$ which is proportional to the left hand side of of (23), see Itin . Then, $$W=(e^\alpha {}_{}{}^{(2)}IH)(e_\alpha F).$$ (24) Since $`IH`$ and $`F`$ are independent, $`W`$ doesn’t vanish in general. For the Maxwell-Lorentz case, $`{}_{}{}^{(2)}IH`$ is zero and, as a consequence, $`W`$ vanishes. Because the electromagnetic skewon energy-momentum has an antisymmetric piece, it would contribute to the first field equation of an Einstein-Cartan-Maxwell (with skewon) system. Thus, we have another (rather indirect) non-Lagrangian type of coupling of the skewon field to gravity. ### 3.3 Energy density The energy density becomes (spatial indices are $`a,b,c,\mathrm{}=1,2,3`$) $`𝒯_0{}_{}{}^{0}\|_{\mathrm{skewon}}^{}`$ $`=`$ $`ϵ^{0klm}F_{0k}\overline{)}S_l{}_{}{}^{n}F_{nm}^{}=ϵ^{0abc}F_{0a}\overline{)}S_b{}_{}{}^{n}F_{nc}^{}`$ (25) $`=`$ $`ϵ^{0abc}\left(F_{0a}\overline{)}S_b{}_{}{}^{0}F_{0c}^{}+F_{0a}\overline{)}S_b{}_{}{}^{d}F_{dc}^{}\right).`$ The first term in the parenthesis vanishes because of its symmetry in $`a`$ and $`c`$. Thus, $`𝒯_0{}_{}{}^{0}\|_{\mathrm{skewon}}^{}`$ $`=`$ $`ϵ^{0abc}F_{a0}\overline{)}S_b{}_{}{}^{d}F_{dc}^{}.`$ (26) Now we can substitute the electric and the magnetic field strengths: $`𝒯_0{}_{}{}^{0}\|_{\mathrm{skewon}}^{}`$ $`=`$ $`E_a\overline{)}S_b{}_{}{}^{d}ϵ_{}^{abc}ϵ_{dce}B^e=E_a\overline{)}S_b{}_{}{}^{a}B_{}^{b}E_a\overline{)}S_b{}_{}{}^{b}B_{}^{a}.`$ We collect the terms and find $$𝒯_0{}_{}{}^{0}\|_{\mathrm{skewon}}^{}=(\overline{)}S_a{}_{}{}^{b}\delta _a^b\overline{)}S_c{}_{}{}^{c})E_bB^a.$$ (27) This is an astonishingly simple premetric result. Note that the second invariant of the electromagnetic field $`I_2:=FF=2d\sigma BE`$ (see , p.126) enters the energy expression inter alia. Recall also that $`\overline{)}S_c{}_{}{}^{c}=\overline{)}S_0^0`$. ### 3.4 Specialization: The spatially isotropic skewon field The spacetime decomposition of the skewon field reads $$\overline{)}S_i{}_{}{}^{j}=\left(\begin{array}{cc}s_c^c& m^a\\ n_b& s_b^a\end{array}\right).$$ (28) Nieves & Pal chose (in nuclear matter) a spatially isotropic skewon field according to $$s_a{}_{}{}^{b}=\frac{s}{2}\delta _a^b,m^a=0,n_a=0(\mathrm{Nieves}\&\mathrm{Pal}).$$ (29) In order to be able to substitute this into (27), we compute $`s_a{}_{}{}^{b}\delta _a^bs_c{}_{}{}^{c}={\displaystyle \frac{s}{2}}\delta _a^b\delta _a^b{\displaystyle \frac{3s}{2}}=s\delta _a^b.`$ (30) We substitute into (27) and find $${}_{}{}^{\mathrm{k}}𝒯_{0}^{}{}_{}{}^{0}\|_{\mathrm{skewon}\mathrm{N}\&\mathrm{P}}^{}=sE_aB^a.$$ (31) Hence the energy density here is proportional to the second invariant $`I_2`$. Since Nieves & Pal didn’t compute the energy of their skewon field, we cannot compare (31) with earlier results. A direct check of (31) starts from the premetric electromagnetic energy density $$u=\frac{1}{2}\left(𝒟^aE_a+_aB^a\right).$$ (32) For the Nieves & Pal skewon we have (, p.262) $$𝒟^a\|_{\mathrm{skewon}\mathrm{N}\&\mathrm{P}}=sB^a,_a\|_{\mathrm{skewon}\mathrm{N}\&\mathrm{P}}=sE_a.$$ (33) Thus, $$u\|_{\mathrm{skewon}\mathrm{N}\&\mathrm{P}}=sE_aB^a,\text{q.e.d.}$$ (34) ## 4 Einstein-Cartan theory with skewon, dilaton, and axion interaction In electrodynamics, one can think of the constitutive 2-form $`𝔎^{\alpha \beta }`$ either as a field determined by the electromagnetic properties of some fixed distribution of background matter or as a property of spacetime itself. Whereas the standard matter fields and the electromagnetic potential are dynamical variables, $`𝔎`$ is a fixed, non-dynamical (or external) field. One can try to describe this situation by an “effective” Lagrangian formalism. Consider, for instance, the simple Lagrangian, quadratic in the field strengths, $`^{}\chi ^{ijkl}F_{ij}F_{kl}`$. It is clear that here only the piece of $`\chi ^{ijkl}`$ symmetric under the exchange $`(i,j)(k,l)`$ survives and thus the related field equations do not contain the skewon piece of $`\chi `$. In a way, this result might have been expected: the skewon field is known to cause dissipative effects in electrodynamics and, consequently, one does not expect to have a simple local Lagrangian description of the complete dynamics. In trying to extend our understanding of $`𝔎`$ to the gravitational sector, we adopt the interpretation of $`𝔎`$ as a property of spacetime, and we will study some of its consequences. ### 4.1 Specialization: principal part with metric and dilaton Although it is known that one can construct a gravitational theory without a metric, all such models are limited to the vacuum case, see, e.g., . It is unclear whether one can construct a viable gravity theory without a metric in the presence of nontrivial matter sources. Accordingly, we will now specialize to the case when the metric field is available as, e.g., in metric-affine gravity (MAG), see . Then, the Maxwell-Lorentz electrodynamics yields the principal part of the form: $${}_{}{}^{(1)}𝔎_{}^{\alpha \beta }=\lambda \eta ^{\alpha \beta }({}_{}{}^{(1)}\kappa _{\gamma \delta }^{}{}_{}{}^{\alpha \beta }=\lambda \widehat{ϵ}_{\gamma \delta \mu \nu }\sqrt{g}g^{\alpha \mu }g^{\beta \nu }).$$ (35) Here $`\eta ^{\alpha \beta }={}_{}{}^{}(\vartheta ^\alpha \vartheta ^\beta )`$ is defined with the help of the Hodge star for the spacetime metric $`g`$, whereas the scalar field $`\lambda (x)`$ is the dilaton field that represents a factor of the principal part of $`𝔎`$ absorbing its physical dimension. The dilaton comes as a companion of the skewon and the axion even on the premetric level. When $`\lambda =`$ const, and the skewon and axion are absent, we recover from (35) the standard Maxwell-Lorentz electrodynamics. Now, we recall that the Einstein-Cartan theory (see Blagojević , Gronwald et al. , and/or Trautman ) is determined by the Lagrangian 4-form (in units $`\kappa =8\pi G=1`$) $$V_{\mathrm{EC}}=\frac{1}{2}\eta ^{\alpha \beta }R_{\alpha \beta }.$$ (36) Noticing that the constitutive 2-form (6) with the principal part (35) provides a natural extension of $`\eta ^{\alpha \beta }`$ by taking into account the electromagnetic companions of the metric, we can propose a generalization of the Einstein-Cartan theory by means of the Lagrangian $`V_{\mathrm{gEC}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝔎^{\alpha \beta }R_{\alpha \beta }={\displaystyle \frac{1}{2}}^{(1)}𝔎^{\alpha \beta }{}_{}{}^{(6)}R_{\alpha \beta }^{}`$ (37) $`+{\displaystyle \frac{1}{2}}^{(2)}𝔎^{\alpha \beta }({}_{}{}^{(2)}R_{\alpha \beta }^{}+{}_{}{}^{(5)}R_{\alpha \beta }^{})+{\displaystyle \frac{1}{2}}^{(3)}𝔎^{\alpha \beta }^{(3)}R_{\alpha \beta }.`$ Here we substituted the irreducible decomposition of the curvature into 6 pieces $`R_{\alpha \beta }=_{A=1}^6{}_{}{}^{(A)}R_{\alpha \beta }^{}`$, see . Since $`{}_{}{}^{(2)}R_{\alpha \beta }^{}`$ is the so-called paircommutator and $`{}_{}{}^{(5)}R_{\alpha \beta }^{}`$ corresponds to the antisymmetric piece of the Ricci tensor, we recognize that in (37) the skewonic part $`{}_{}{}^{(2)}𝔎_{}^{\alpha \beta }`$ couples only to these specific post-Riemannian pieces of the curvature. More generally, the contributions of the skewon and axion are only nontrivial for a Riemann-Cartan geometry with a nonvanishing torsion 2-form $`T^\alpha `$. The torsion itself can be also irreducibly decomposed according to $`T^\alpha ={}_{}{}^{(1)}T_{}^{\alpha }+{}_{}{}^{(2)}T_{}^{\alpha }+{}_{}{}^{(3)}T_{}^{\alpha }`$, with the second and the third irreducible torsion pieces defined as usual by $${}_{}{}^{(2)}T_{}^{\alpha }=\frac{1}{3}\vartheta ^\alpha T,{}_{}{}^{(3)}T_{}^{\alpha }=\frac{1}{3}{}_{}{}^{}(\vartheta ^\alpha P).$$ (38) The 1-forms of the trace and the axial trace of torsion are introduced by $`T:=e_\alpha T^\alpha `$ and $`P:={}_{}{}^{}(\vartheta ^\alpha T_\alpha )`$, respectively. By making use of the first Bianchi identity $`DT^\alpha =R_\beta {}_{}{}^{\alpha }\vartheta ^\beta `$, we can rewrite the above Lagrangian (37) into an equivalent form $$V_{\mathrm{gEC}}=\frac{1}{2}\lambda \eta ^{\alpha \beta }{}_{}{}^{(6)}R_{\alpha \beta }^{}+(D\overline{)}S^\alpha )({}_{}{}^{(1)}T_{\alpha }^{}+{}_{}{}^{(2)}T_{\alpha }^{})\frac{1}{2}D(\alpha \vartheta ^\alpha ){}_{}{}^{(3)}T_{\alpha }^{}+d\mathrm{\Psi }.$$ (39) Here the total derivative term contains the 3-form $`\mathrm{\Psi }:=\overline{)}S^\alpha ({}_{}{}^{(1)}T_{\alpha }^{}+{}_{}{}^{(2)}T_{\alpha }^{})+\frac{1}{2}\alpha \vartheta ^\alpha {}_{}{}^{(3)}T_{\alpha }^{}`$. Obviously, the skewon field $`\overline{)}S^\alpha `$ couples to the tensor and the vector pieces of the torsion, the axion field $`\alpha `$, however, to the axial torsion (totally antisymmetric torsion). ### 4.2 Generalized gravitational field equations The general framework for the derivation of the field equations is provided by the Noether-Lagrange machinery developed in the review paper , see its Sec. 5.8.1. The gravitational field equations are given by the system of the so-called first and the second field equations of gravity: $`DH_\alpha E_\alpha `$ $`=`$ $`\mathrm{\Sigma }_\alpha ,`$ (40) $`DH_{\alpha \beta }+\vartheta _{[\alpha }H_{\beta ]}`$ $`=`$ $`\tau _{\alpha \beta }.`$ (41) The sources arise as the variational derivatives of the material Lagrangian with respect to the coframe and the connection, and they represent the canonical energy-momentum current $`\mathrm{\Sigma }_\alpha `$ and the spin current $`\tau _{\alpha \beta }`$, respectively. For the Lagrangian (37) of the generalized Einstein-Cartan theory we find straightforwardly the gravitational gauge field momenta $$H_\alpha =\frac{V_{\mathrm{gEC}}}{T^\alpha }=0,H_{\alpha \beta }=\frac{V_{\mathrm{gEC}}}{R^{\alpha \beta }}=\frac{1}{2}𝔎_{\alpha \beta },$$ (42) and the gravitational canonical energy 3-form $$E_\alpha =\frac{1}{2}(e_\alpha 𝔎_{\beta \gamma })R^{\beta \gamma }.$$ (43) Accordingly, the gravitational field equations read $`{\displaystyle \frac{1}{2}}(e_\alpha 𝔎_{\beta \gamma })R^{\beta \gamma }`$ $`=`$ $`\mathrm{\Sigma }_\alpha ,`$ (44) $`{\displaystyle \frac{1}{2}}D𝔎_{\alpha \beta }`$ $`=`$ $`\tau _{\alpha \beta }.`$ (45) Explicitly, the first equation has the form $$\frac{1}{2}\lambda \eta _{\alpha \beta \gamma }R^{\beta \gamma }\overline{)}S^\beta R_{\beta \alpha }+(e_\alpha \overline{)}S^\beta )\vartheta ^\gamma R_{\beta \gamma }\alpha \vartheta ^\beta R_{\alpha \beta }=\mathrm{\Sigma }_\alpha .$$ (46) Besides the first Einsteinian term (modified by the scalar dilaton coupling à la Brans-Dicke), we now see that the skewon and the axion fields bring into the first equation new terms which all depend on the Riemann-Cartan curvature. The second field equation determines the spacetime torsion in terms of the spin current of matter and the additional contributions of skewon, axion, and dilaton. Explicitly, we have: $$\frac{\lambda }{2}\eta _{\alpha \beta \gamma }T^\gamma \frac{1}{2}\eta _{\alpha \beta }d\lambda +T_{[\alpha }\overline{)}S_{\beta ]}\vartheta _{[\alpha }D\overline{)}S_{\beta ]}\frac{1}{2}\vartheta _\alpha \vartheta _\beta d\alpha \alpha T_{[\alpha }\vartheta _{\beta ]}=\tau _{\alpha \beta }.$$ (47) In principle, one can resolve this algebraic equation with respect to the components of the torsion. In vacuum, when the matter spin vanishes, we find that the torsion is determined exclusively by the metric companion fields: the dilaton, the skewon, and the axion (and their derivatives). There is a particular exact solution of the field equations in vacuum, which corresponds to a teleparallel geometry. Indeed, for $`\mathrm{\Sigma }=\tau =0`$, we see that $`R_{\alpha \beta }=0`$ solves the first field equation, while (47) defines the intrinsic torsion of spacetime in terms of its $`𝔎`$-structure (dilaton, skewon, and axion). ### 4.3 Simple vacuum solution Unfortunately, although the second field equation looks rather simple, it is not easy to find the torsion components from it explicitly. Nevertheless, we can illustrate how the theory works in a simple case when the skewon is absent. Then, the terms with $`\overline{)}S`$ disappear from the equation, and we find that the vacuum torsion has the following simple form, $$T^\alpha ={}_{}{}^{(2)}T_{}^{\alpha }+{}_{}{}^{(3)}T_{}^{\alpha },$$ (48) that is, the tensor piece $`{}_{}{}^{(1)}T_{}^{\alpha }`$ of the torsion vanishes. From (47), we find: $`T`$ $`=`$ $`{\displaystyle \frac{3/2}{\lambda ^2+\alpha ^2}}(\lambda d\lambda +\alpha d\alpha ),`$ (49) $`P`$ $`=`$ $`{\displaystyle \frac{3}{\lambda ^2+\alpha ^2}}(\alpha d\lambda \lambda d\alpha ).`$ (50) This can be verified if we notice that the following identities hold true in exterior calculus: $`\eta _{\alpha \beta \gamma }{}_{}{}^{(2)}T_{}^{\gamma }={\displaystyle \frac{2}{3}}\eta _{\alpha \beta }T,`$ $`\eta _{\alpha \beta \gamma }{}_{}{}^{(2)}T_{}^{\gamma }={\displaystyle \frac{1}{3}}\vartheta _\alpha \vartheta _\beta P,`$ (51) $`{}_{}{}^{(2)}T_{[\alpha }^{}\vartheta _{\beta ]}={\displaystyle \frac{1}{3}}\vartheta _\alpha \vartheta _\beta T,`$ $`{}_{}{}^{(3)}T_{[\alpha }^{}\vartheta _{\beta ]}={\displaystyle \frac{1}{6}}\eta _{\alpha \beta }P.`$ (52) We thus conclude that the trace and the axial trace 1-forms of the torsion are determined, in vacuum, by the dilaton and the axion fields. In particular, when the axion is absent, $`\alpha =0`$, we recover a Brans-Dicke type gravity constructed in the Einstein-Cartan framework in . In that case the axial trace vanishes, whereas the torsion trace is proportional to the gradient of the dilaton field. Otherwise, for the case of a constant dilaton, $`\lambda =\lambda _0=`$const, both torsion 1-forms are nontrivial and depend on the axion only. Now, when we return to the general case, it is straightforward to verify that the nontrivial skewon induces the trace-free irreducible part of torsion, $`{}_{}{}^{(1)}T_{}^{\alpha }`$, in addition to the trace and the axial trace 1-forms. ## 5 Gravitational energy in the generalized Einstein-Cartan theory We now wish to study the energy of the generalized Einstein-Cartan theory (37) in the framework of the canonical formalism . Using simple algebra, we can rewrite the gravitational piece of (37) in the form $`V_{\mathrm{gEC}}=_{\mathrm{gEC}}d^4x`$, with $$_{\mathrm{gEC}}=\frac{1}{4}\chi _{\alpha \beta }{}_{}{}^{ij}R_{ij}^{}{}_{}{}^{\alpha \beta }(\mathrm{\Gamma }).$$ (53) Here, $`\chi _{\alpha \beta }{}_{}{}^{ij}=2\lambda \sqrt{g}e^i{}_{[\alpha }{}^{}e_{}^{j}{}_{\beta ]}{}^{}+\overline{\chi }_{\alpha \beta }{}_{}{}^{ij},`$ (54) $`\overline{\chi }_{\alpha \beta }{}_{}{}^{ij}=e_{k\alpha }e_{l\beta }\left(2ϵ^{klm[i}\mathit{}_m{}_{}{}^{j]}+\alpha ϵ^{klij}\right).`$ (55) The Latin and Greek indices are raised and lowered with the help of the spacetime metric $`g_{ij}=e_i{}_{}{}^{\alpha }e_{j}^{}{}_{}{}^{\beta }g_{\alpha \beta }^{}`$ and the Lorentz metric $`g_{\alpha \beta }=\mathrm{diag}(+1,1,1,1)`$, respectively. Primary constraints are similar to those of the standard Einstein-Cartan theory: $`\pi _\alpha {}_{}{}^{0}0,\pi _{\alpha \beta }{}_{}{}^{0}0,`$ $`\pi _\alpha {}_{}{}^{a}0,\varphi _{\alpha \beta }{}_{}{}^{a}:=\pi _{\alpha \beta }{}_{}{}^{a}\chi _{\alpha \beta }{}_{}{}^{0a}0.`$ Since the Lagrangian is linear in the velocities $`\dot{\mathrm{\Gamma }}`$, the canonical Hamiltonian is given by $`_c=_{\mathrm{gEC}}(\dot{\mathrm{\Gamma }}=0)`$: $$_c=\frac{1}{4}\chi _{\alpha \beta }{}_{}{}^{ab}R_{ab}^{}{}_{}{}^{\alpha \beta }+\frac{1}{2}\mathrm{\Gamma }_0{}_{}{}^{\alpha \beta }_{a}^{}\chi _{\alpha \beta }{}_{}{}^{0a}+_aU^a,$$ (56) where $`U^a=\chi _{\alpha \beta }{}_{}{}^{0a}\mathrm{\Gamma }_{0}^{}{}_{}{}^{\alpha \beta }/2`$. Looking at the form of $`\chi _{\alpha \beta }^{ij}`$, we see that the canonical Hamiltonian contains not only the standard Einstein-Cartan piece, modified by the presence of dilaton, but also an additional skewon-axion contribution, $$_c^{\mathrm{SA}}=\frac{1}{4}\overline{\chi }_{\alpha \beta }{}_{}{}^{ab}R_{ab}^{}{}_{}{}^{\alpha \beta }+\frac{1}{2}\mathrm{\Gamma }_0{}_{}{}^{\alpha \beta }_{a}^{}\overline{\chi }_{\alpha \beta }{}_{}{}^{a0}+_aU_{\mathrm{SA}}^a,$$ (57) with $`U_{\mathrm{SA}}^a=\overline{\chi }_{\alpha \beta }{}_{}{}^{0a}\mathrm{\Gamma }_{0}^{}{}_{}{}^{\alpha \beta }/2`$. The total Hamiltonian has the form $$_T=_c+u^\alpha {}_{0}{}^{}\pi _{\alpha }^{}{}_{}{}^{0}+\frac{1}{2}u^{\alpha \beta }{}_{0}{}^{}\pi _{\alpha \beta }^{}{}_{}{}^{0}+u^\alpha {}_{a}{}^{}\pi _{\alpha }^{}{}_{}{}^{a}+\frac{1}{2}u^{\alpha \beta }{}_{a}{}^{}\varphi _{\alpha \beta }^{}{}_{}{}^{a}.$$ (58) The simple Hamiltonian structure obtained so far is sufficient to derive the canonical expression for the gravitational energy. In the Hamiltonian formalism, symmetry properties of a dynamical system are described by the canonical generators $`G[\phi ,\pi ]`$. Since they act on basic dynamical variables via Poisson brackets, they must be differentiable. A local functional $`F[\phi ,\pi ]=d^3xf(\phi ,\phi ,\pi ,\pi )`$ has well defined functional derivatives if its variation can be written in the form $`\delta F[\phi ,\pi ]=d^3x(A\delta \phi +B\delta \pi )`$, where terms of the form $`\delta (\phi )`$ and $`\delta (\pi )`$ are absent. If the generator $`G[\phi ,\pi ]`$ is not differentiable, its form can be improved by adding a suitable surface term, whereupon it becomes differentiable. On shell, these surface terms represent the values of the related conserved charges. The canonical generator of time translations is defined by the total Hamiltonian: $$P_0=d^3x\widehat{}_T,_T\widehat{}_T+_\alpha U^\alpha .$$ (59) Looking at the skewon-axion piece of $`P_0`$, we find that its variation has the form $$\delta P_0^{\mathrm{SA}}=d^3x\delta \widehat{}_c^{\mathrm{SA}}+N=\frac{1}{2}d^3x_a(\overline{\chi }_{\alpha \beta }{}_{}{}^{ab}\delta \mathrm{\Gamma }_b{}_{}{}^{\alpha \beta })+N,$$ where $`N`$ denotes well defined, normal (regular) terms. Thus, $`P_0^{\mathrm{SA}}`$ can be made differentiable by adding a surface term: $`P_0^{\mathrm{SA}}\stackrel{~}{P}_0^{\mathrm{SA}}=P_0^{\mathrm{SA}}+^{\mathrm{SA}},`$ $`^{\mathrm{SA}}={\displaystyle \frac{1}{2}}{\displaystyle }dS_a(\overline{\chi }_{\alpha \beta }{}_{}{}^{ab}\mathrm{\Gamma }_{b}^{}{}_{}{}^{\alpha \beta }).`$ (60) In order to ensure the convergence of the surface integral $`^{\mathrm{SA}}`$, we have to adopt suitable asymptotic conditions. For localized gravitational sources (matter fields decrease sufficiently fast at large distances and give no contribution to surface integrals), we can assume that spacetime is asymptotically flat. The related asymptotic conditions for the gravitational variables, when expressed in the standard spherical coordinates, take the simple form: $$e_i{}_{}{}^{\alpha }=\delta _i^\alpha +𝒪(1/r),\mathrm{\Gamma }_i{}_{}{}^{\alpha \beta }=𝒪(1/r^2).$$ (61) Hence, $`^{\mathrm{SA}}`$ is convergent if $$\overline{\chi }_{\alpha \beta }{}_{}{}^{ab}\mathrm{const}.\mathrm{for}r\mathrm{}.$$ (62) The surface term $`^{\mathrm{SA}}`$ represents the value of the skewon-axion contribution to the gravitational energy. It is produced by the interaction between the skewon-axion term $`\overline{\chi }`$, and the connection $`\mathrm{\Gamma }`$. One should remember that the complete gravitational energy contains also the standard Einstein-Cartan piece, modified by the presence of the dilaton. The adopted asymptotics, extended naturally to the dilaton field by $$\lambda (x)\mathrm{const}.\mathrm{for}r\mathrm{},$$ ensures the conservation of the gravitational energy. ## 6 Concluding remarks (1) According to its definition, the skewon field is some kind of permeability/permittivity of spacetime — and this in a premetric setting when the metric has not yet ”condensed”. In this sense, the skewon field is an elementary electromagnetic property of spacetime. As such, it influences light propagation. (2) The skewon field contributes non-trivially to the electromagnetic energy. In particular, it induces an asymmetric electromagnetic energy-momentum tensor, which can cause specific gravitational effects as a source term in the Einstein-Cartan-Maxwell system (with skewon). (3) A smooth deformation of the Einstein-Cartan theory has been introduced and studied as a simple dynamical model incorporating gravitational effects of the skewon field. We found the generalized gravitational field equations and were able to determine the contribution of the skewon field to the gravitational energy. ### Acknowledgments We are very grateful to Yakov Itin (Jerusalem) and to Christian Heinicke (Cologne) for various helpful remarks. G. Rubilar would like to thank the Fundacion Andes (Convenio C-13860) for financial support. ## Appendix. Decomposition of the local and linear constitutive law We start with a local (in space and in time) and linear constitutive law $$IH=\kappa [F].$$ (63) The operator $`\kappa `$ acts on the electromagnetic field strength $`F`$. Because of its linearity, we have for any 2-forms $`\mathrm{\Psi },\mathrm{\Phi }`$, if $`a,b`$ are two real scalar factors, $$\kappa [a\mathrm{\Psi }+b\mathrm{\Phi }]=a\kappa [\mathrm{\Psi }]+b\kappa [\mathrm{\Phi }].$$ (64) We substitute the decomposition of the field strength $`F`$ into (63): $$IH=\kappa [\frac{1}{2}F_{\alpha \beta }\vartheta ^\alpha \vartheta ^\beta ]=\frac{1}{2}\kappa [\vartheta ^\alpha \vartheta ^\beta ]F_{\alpha \beta }.$$ (65) We introduce the constitutive tensor-valued 2-form $$𝔎^{\alpha \beta }:=\kappa [\vartheta ^\alpha \vartheta ^\beta ]$$ (66) and recall the decomposition $`F_{\alpha \beta }=e_\beta e_\alpha F`$. Then the constitutive relation can be brought into the compact form $$IH=\frac{1}{2}𝔎^{\alpha \beta }e_\beta e_\alpha F,\text{with}𝔎^{\alpha \beta }=𝔎^{\beta \alpha }.$$ (67) Here the 2-form $`𝔎^{\alpha \beta }`$ decomposes as $$𝔎^{\alpha \beta }=\frac{1}{2}\kappa _{\gamma \delta }{}_{}{}^{\alpha \beta }\vartheta _{}^{\gamma }\vartheta ^\delta .$$ (68) Contractions yield a 1-form and a 0-form, respectively: $$𝔎^\beta :=e_\alpha 𝔎^{\alpha \beta }=\kappa _{\alpha \delta }{}_{}{}^{\alpha \beta }\vartheta _{}^{\delta }=:\kappa _\delta {}_{}{}^{\beta }\vartheta _{}^{\delta },𝔎:=e_\beta 𝔎^\beta =\kappa _\beta ^\beta =:\kappa .$$ (69) The tracefree part of the 1-form is $$\overline{)}𝔎^\alpha :=𝔎^\alpha \frac{1}{4}𝔎\vartheta ^\alpha .$$ (70) In this way we can decompose the constitutive antisymmetric tensor valued 2-form into its 3 irreducible pieces, $$𝔎^{\alpha \beta }=^{(1)}𝔎^{\alpha \beta }+^{(2)}𝔎^{\alpha \beta }+^{(3)}𝔎^{\alpha \beta },$$ (71) with the skewon and the axion pieces $$^{(2)}𝔎^{\alpha \beta }:=\overline{)}𝔎^{[\alpha }\vartheta ^{\beta ]}\text{and}^{(3)}𝔎^{\alpha \beta }:=\frac{1}{12}𝔎\vartheta ^\alpha \vartheta ^\beta .$$ (72) The factors can be determined with some trivial algebra. Note the constraints $$e_\alpha ^{(1)}𝔎^{\alpha \beta }=0\mathrm{and}e_\alpha \overline{)}𝔎^\alpha =0.$$ (73) The irreducible pieces in (72) can also be written in components. With the help of the generalized Kronecker delta (see ), we find $$^{(2)}\kappa _{\gamma \delta }{}_{}{}^{\alpha \beta }=2\delta _{\gamma \delta }^{\epsilon [\alpha }\overline{)}S_\epsilon {}_{}{}^{\beta ]},_{}^{(3)}\kappa _{\gamma \delta }{}_{}{}^{\alpha \beta }=\delta _{\gamma \delta }^{\alpha \beta }\alpha .$$ (74) If we substitute (71) into (67) and observe $`\vartheta ^\alpha (e_\alpha \omega )=p\omega `$, where $`\omega `$ is a $`p`$-form, then we finally have $`IH`$ $`=`$ $`{}_{}{}^{(1)}IH+{}_{}{}^{(2)}IH+{}_{}{}^{(3)}IH`$ (75) $`=`$ $`{\displaystyle \frac{1}{2}}({}_{}{}^{(1)}𝔎_{}^{\alpha \beta }e_\beta e_\alpha +\overline{)}𝔎^\alpha e_\alpha +{\displaystyle \frac{1}{6}}𝔎)F.`$ Thus, the principal part of the constitutive 2-form $`𝔎^{\alpha \beta }`$ is represented by the $`[_0^2]`$ antisymmetric tensor-valued 2-form $`{}_{}{}^{(1)}𝔎_{}^{\alpha \beta }=^{(1)}𝔎^{\beta \alpha }`$, the skewon part by the vector-valued 1-form $`\overline{)}𝔎^\alpha `$, and the axion part by the pseudoscalar $`𝔎`$. The translation into our usual language is made by $`\overline{)}S^\alpha =\frac{1}{2}\overline{)}𝔎^\alpha `$ and $`\alpha =\frac{1}{12}𝔎.`$ Incidentally, the IB-medium of Lindell is defined by $`{}_{}{}^{(1)}𝔎_{}^{\alpha \beta }=0`$. If additionally $`{}_{}{}^{(2)}𝔎_{}^{\alpha \beta }=0`$ (vanishing skewon field), then only $`{}_{}{}^{(3)}𝔎_{}^{\alpha \beta }=\frac{1}{12}𝔎\vartheta ^\alpha \vartheta ^\beta `$ is left over, the axion field with 1 component, or, in the language of Lindell & Sihvola , the perfect electromagnetic conductor (PEMC). =========
warning/0506/astro-ph0506530.html	ar5iv	text	# Spectral Indications of thermal Sunyaev-Zel’dovich Effect in Archeops and WMAP Data ## 1 Introduction The thermal Sunyaev-Zel’dovich effect (hereafter tSZ, Sunyaev & Zel’dovich 1980 ) constitutes a unique tool to explore the presence of baryons in the Universe. It arises as a consequence of the distortion that the black body spectrum of the Cosmic Microwave Background (CMB) radiation experiences when it encounters a hot electron plasma. In this Compton scattering, electrons transfer energy to the CMB radiation, generating an excess of high energy photons and a deficit in the low energy tail of the distribution. This photon reallocation translates into a frequency dependent change of the brightness temperature of the CMB, which, in the non-relativistic limit, has a very simple form, ($`f_{tSZ}(x)=x/\mathrm{tanh}[x/2]4`$, with $`x=h\nu /k_BT_{CMB}`$ the adimensional frequency in terms of the CMB temperature monopole $`T_{CMB}`$). The amplitude of this distortion is proportional to the electron pressure integrated along the line of sight ($`\delta T_{tSZ}/T_{CMB}=f_{tSZ}(x)𝑑r\sigma _Tn_ek_BT_e/(m_ec^2)`$, with $`\sigma _T`$ the Thomson cross section and $`n_e`$, $`T_e`$ and $`m_e`$ the electron number density, temperature and mass, respectively); and this makes this effect particularly sensitive to collapsed or collapsing structures containing hot electrons, such as clusters and superclusters of galaxies (see Birkinshaw 1999 for a extensive review). In addition to the intrinsic energy inhomogeneities of CMB radiation generated during inflation, there are further temperature anisotropies introduced in the CMB during recombination, which are mainly caused by two physical processes. These processes are the last Doppler kick exerted by electrons via Thomson scattering just before recombining, and the subsequent gravitational redshift experienced by CMB photons as they climb the potential wells generated by the inhomogeneities in the matter distribution (Sachs & Wolfe effect, e.g. Hu & Sugiyama 1995 ). While all this happens at $`z1100`$, a similar scenario can take place at much lower redshifts: as the first stars reionise the universe, new free electrons are produced which again scatter CMB photons, partially blurring primordial anisotropies generated during recombination and introducing new ones at much larger angular scales. Also, if $`\mathrm{\Omega }_\mathrm{\Lambda }`$ is non-zero, the decay of gravitational potentials in linear scales introduces a net blueshift in the CMB radiation at late epochs (z $`<`$ 2), which is known as the Integrated Sachs Wolfe effect, (ISW Sachs & Wolfe 1967 ). Despite the fact that most of these phenomena introduce temperature fluctuations of amplitudes larger than the tSZ effect, the particular frequency dependence of the latter should enable its separation. Whereas the first generation of CMB experiments like COBE Smoot et al.1991 and Tenerife Gutiérrez et al. 2000 aimed to simply detect the largest CMB temperature anisotropies in the big angular scales, experiments like, e.g., Boomerang Mauskopf et al.2000 , VSA Rubiño-Martín et al. 2003 , Archeops Benoît et al. 2002 and WMAP Bennett et al. 2003 have already reached the sensitivity and angular resolution levels required to probe relatively weak signals like the tSZ effect. In this work we perform a combined analysis of Archeops and WMAP CMB data, searching for spectral signatures of the tSZ effect. Previous works Bennett et al. 2003 ; Hernández-Monteagudo & Rubiño-Martín 2004 ; Myers et al.2004 ; Hernández-Monteagudo et al. 2004 ; Afshordi et al. (2004); Fosalaba et al. 2003 ; Fosalaba & Gaztanaga 2004 have claimed the detection of tSZ in WMAP data at different significance levels. However, all those studies were exclusively based on spatial cross-correlations of large scale structure catalogues with CMB data. In this work, we take advantage of the frequency coverage provided by the combination of Archeops and WMAP experiments in order to include an analysis of the frequency behavior of a signal which is spatially correlated with regions hosting large galaxy overdensities. The sketch of the paper is as follows: in Section 2 we summarize the outcome of the Archeops and WMAP experiments, whose data products are analyzed as explained in Section 3. Section 4 shows our results, which are compared with those obtained from WMAP data. Their implications are discussed in Section 5. Finally, we conclude in Section 6. ## 2 The Archeops and WMAP data set The Archeops Benoît et al. 2002 <sup>1</sup><sup>1</sup>1see http://www.archeops.org experiment was designed to obtain a large sky coverage of CMB temperature anisotropies in a single balloon flight at millimeter and submillimeter wavelengths. Archeops is a precursor to the Planck HFI instrument Lamarre et al. 2003 , using the same optical design and the same technology for the detectors, spider–web bolometers, and their cooling 0.1 K dilution fridge. The instrument consists of a 1.5 m aperture diameter telescope and an array of 21 photometric pixels operating at 4 frequency bands centered at 143, 217, 353 and 545 GHz. The two low frequencies are dedicated to CMB studies while high frequency bands are sensitive to foregrounds, essentially to interstellar dust and atmospheric emission. Observations are carried out by spinning the payload around its vertical axis at 2 rpm. Thus the telescope produces circular scans at a fixed elevation of $`41`$ deg. The data were taken during the Arctic night of February 7, 2002 after the instrument was launched by CNES from the Esrange base near Kiruna (Sweden). The entire data set covers $`30`$% of the sky in 12 hours of night observations. For the purpose of this paper, we concentrate in the low frequency channels at 143 and 217 GHz. Maps for each of the bolometers have been produced from the Archeops processed and foreground cleaned timelines, using the Mirage optimal map making code Yvon & Mayet 2005 as discussed in Tristram et al. 2005 . The maps for the 4 most sensitive bolometers at 143 GHz were combined into a single map at 143 GHz and equally the two most sensitive bolometer maps were combined at 217 GHz. The CMB dipole is the prime calibrator of the instrument. The absolute calibration error against the dipole as measured by COBE/DMR Fixsen et al. 1994 and confirmed by WMAP Bennett et al. 2003 is estimated to be 4% and 8% in temperature at 143 GHz and 217 GHz respectively. These errors are dominated by systematic effects. The noise contribution in the combined maps at 143 and 217 GHz was computed using Monte Carlo simulations. For each bolometer, by using the power spectrum of the noise in the time domain data set, we produced fake timelines of Archeops noise. These were processed and projected into maps following the same procedures used for the Archeops data themselves as described before. The WMAP satellite mission was designed to measure the CMB temperature and polarization anisotropies in 5 frequency bands, 23, 33, 41, 61 and 94 GHz with a full sky coverage. The satellite was launched in June 2001 and its first results after the first year of observations Bennett et al. 2003 included the CMB temperature and temperature-polarization cross-correlation power spectra, as well as full sky temperature maps for each of the frequency bands. In this paper we consider only data from the high frequency channels, Q \[41 GHz\], V\[61 GHz\] and W \[94 GHz\], since only for these bands there are foreground clean maps available at the LAMBDA site http://lambda.gsfc.nasa.gov/. ## 3 The Statistical Analysis In this Section we outline the correlation analysis performed on Archeops and WMAP data and the 2MASS Extended Source Catalog (XSC, Jarrett et al. 2003 ) on the sky region covered by Archeops. This analysis is essentially identical to that applied in Hernández-Monteagudo et al. 2004 , paper to which we refer for a more detailed description of the statistical method. It consists of a pixel-to-pixel comparison of Archeops CMB data with a template of the large scale structure built from the 2MASS XSC catalog. The 2MASS XSC catalog contains approximately 1.6 million galaxies detected in the infrared filters I, J and K, and covers the whole celestial sphere. Those frequencies are particularly insensitive to dust absorption, and for this reason this catalog can trace the extragalactic structure at very low galactic latitudes. The galaxy templates built from it take into account the spatial distribution of the galaxies and the instrumental beam of the CMB experiment. By using the HEALPix<sup>2</sup><sup>2</sup>2HEALPix’s URL site: http://www.eso.org/science/healpix/ Górski et al. 1999 tessellation of the sky, we built a map with the same resolution parameter ($`N_{side}=512`$) than the one used in Archeops and WMAP data. Every pixel was assigned a value equal to the number of galaxies present in such pixel. The resulting map was then convolved with a window function corresponding to the instrumental beam of each of the detectors taken into consideration. For Archeops, the resulting templates were then weighted by their noise levels and co-added per frequency band, in such a way that we ended up with two different galaxy templates corresponding to the 143 GHz and 217 GHz bands. For WMAP, we produced templates for the Q, V and W bands. In the next step, we sorted the pixels of each template in terms of its amplitude, so that first pixels would have higher galaxy densities. These pixels were grouped in patches of varying sizes (32, 64 or 128 pixels per patch), and again patches were sorted in such a way that first patches contained larger projected galaxy densities. Next, we analyzed each of these patches separately, by comparing them to the corresponding patches in the Archeops and WMAP CMB maps on a pixel-to-pixel basis. As explained in e.g. Hernández-Monteagudo & Rubiño-Martín 2004 , it is possible to estimate the contribution of a given spatial template (M) on a total measured temperature map (T), which is the result of the addition of several components: $$\text{T}=\text{T}_{cmb}+\stackrel{~}{\alpha }\text{M}+\text{N},$$ (1) namely CMB ($`\text{T}_{cmb}`$), instrumental noise (N) and some signal coming from the extragalactic template (M). The contribution of M is parametrized by $`\stackrel{~}{\alpha }`$, and an optimal value of it (optimal in terms of the temperature model given in eq.(1)), together with its formal error bar, is given by $$\alpha =\frac{\text{T}𝒞^1\text{M}^T}{\text{M}𝒞^1\text{M}^T},\sigma _\alpha =\sqrt{\frac{1}{\text{M}𝒞^1\text{M}^T}}.$$ (2) The matrix $`𝒞`$ is the covariance matrix of T, which must contain the correlation matrices of both the CMB and noise components. For the small scales we are probing here (the pixel has a size of $`7`$ arcmins), the noise is the main contributor to the covariance matrix<sup>3</sup><sup>3</sup>3Note that for these small scales, the assumption of a piossonian noise is a good approximation. In our case, this matrix must be evaluated only for the pixels belonging to the patch under analysis. Since our patches are relatively small, the inversion of this matrix poses no numerical problem. Therefore, for every patch a value of $`\alpha `$ and $`\sigma _\alpha `$ was obtained. However, uncertainties in the determination of the noise amplitude may bias our determination of $`\sigma _\alpha `$, and for this reason we computed a different estimate of the uncertainty of $`\alpha `$, namely the r.m.s variation of this parameter for all available patches, which will be denoted as $`\sigma _\alpha ^{rms}`$. We remark the fact that, according to eq.2, an error in the noise normalization does not affect the estimates of $`\alpha `$, but only those of $`\sigma _\alpha `$. When comparing $`\sigma _\alpha `$ with $`\sigma _\alpha ^{rms}`$, we found that for the 143 GHz channel of Archeops, the latter was about a 40% larger than the former, and hence we decided to adopt it in order to quote conservative estimates of statistical significance. For the 217 GHz case no such bias was found, and we decided to use again $`\sigma _\alpha ^{rms}`$. ## 4 Combined results for Archeops and WMAP In figure 1 we plot the recovered $`\alpha `$’s versus the patch index. Results are grouped in six different panels: top and bottom panels refer to WMAP and Archeops experiments, respectively, whereas left, middle and right panels display results for patches of 32, 64 and 128 pixels, respectively. For Archeops, filled circles and diamonds refer to 143 GHz and 217 GHz respectively, whereas for WMAP those symbols correspond to the W and the Q channels, being the results of the V band given by the crosses. Dark and light grey colored bands limit the 2-$`\sigma `$ confidence levels for filled circles and diamonds, respectively, while the moderately dark bands refer to the V band (crosses) in the case of WMAP. As explained above, for the 143 GHz and 217 GHz channels the amplitude of the shaded regions was computed from the typical dispersion of the values obtained for $`\alpha `$ in patches where the tSZ contribution is expected to be negligible, i.e., in patches with indexes between 40 and 300. We have found that for the 143 GHz channel the first patch contains an unusual negative $`\alpha `$ while at 217 GHz seems to be compatible with zero. This patch, in the case it contains 64 pixels, hosts the central pixels of 20 different ACO clusters of galaxies, COMA among them. Out of them, 11 (COMA again included) are already sampled by the 32 densest pixels. Its statistical significance is slightly bigger for patches with 64 pixels ($`>`$ 2.5-$`\sigma ^{rms}`$), since it contains the first two very negative patches of 32 pixels each, (see figure 1d). In figure 1e, from the first 300 patches, very few of them ($`12`$) depart from zero by an amount similar to that of the first patch; and such number is very close to what one expects under Gaussian statistics at 2.5-$`\sigma `$ level of significance. This peculiar behavior of the first (or first two) patch(es) disappears at 217 GHz (see diamonds): in no case the diamonds corresponding to the first two patches trespass the 2-$`\sigma ^{rms}`$ limit. This picture is consistent with part of the signal being generated by the tSZ effect, since such component is negative at 143 GHz and becomes zero at 217 GHz. In order to interpret the results from WMAP, one musts keep in mind that the Q and V bands have remarkably larger beams than the W band: while the Archeops’ and W band’s Point Spread Functions are similar in size ($`13`$ arcmins), the beams of the Q and V bands have an average (linear) size of $``$ 31 and 21 arcmin respectively. This, in terms of tSZ flux, corresponds to factors $``$ 5.7 and 2.6 smaller in the Q and V band for point-like objects. On the other hand, it is clear from figure 1 that WMAP has a much lower noise level when compared to Archeops, (approx. a factor of 3.5). As one looks at the top panels in figure 1, one finds that for the first patch of both 32 and 64 pixel size, the W band gives a decrement about 4-$`\sigma `$ away from zero, (which however is not found for the second patch of 32 pixel size). This statistical singularity of the first patch decreases remarkably in the Q and V bands, but still remains close to the 2-$`\sigma `$ level. However, due to the argument on beam dilution on bands Q and V with respect to W, we still find the amplitudes given by the three bands of WMAP consistent with being (at least partially) generated by unresolved objects causing a temperature decrement. It also worth to remark that, if the first patch contains pixels which do not contribute considerably to the tSZ decrement, the significance of the overall $`\alpha `$ obtained for that patch will diminish accordingly. This is our explanation for the decreasing significance of the values of $`\alpha `$ in the first patch for a size of 128 pixels when compared to a size of 64 pixels. ## 5 Discussion ### 5.1 Spectral dependence of the cross-correlation coefficients As discussed in the previous section, Fig. 1 indicates a significant temperature decrement in patch 1 for both the Archeops and WMAP data sets which seems to be compatible with the tSZ effect. For a better assessment of this result we have compared, via a linear fit, the observed spectral dependency of the correlation coefficients $`\alpha `$ in each of the patches to the following model $$\alpha (\nu )=A\times f_{tSZ}(\nu )+n(\nu )$$ (3) where $`A`$ is global calibration factor which is estimated from the fit, $`f_{tSZ}(\nu )`$ represents the spectral behavior of the tSZ-induced change in brightness temperature and $`n(\nu )`$ is the instrumental noise in $`\alpha (\nu )`$. For this model if the signal observed is compatible with the tSZ effect we expect a temperature decrement $`A`$ to be significantly positive, otherwise we expect no correlation and therefore $`A`$ to be compatible with zero. We must remark that, due to the different size of the beams for every channel, we had to rescale the $`\alpha `$’s to a common reference beam size of 13 arcmin. When doing this, we assumed that the signal was generated by unresolved sources, and hence we scaled the $`\alpha `$’s as the ratio of the area of the beam with the reference one. The main results of the fitting procedure described before are shown on figure 2 where we trace the fitted amplitude, $`A`$, of the tSZ signal as a function of the patch number for patches of 32 pixels. We observe that only for patch number 1, where most massive clusters might be, there is a significant temperature decrement ( $`A>0`$ ). For a closer look to the fits, figures 3, 4, 5 represent the $`\alpha `$ coefficient in CMB temperature units as function of frequency in GHz for the 32, 64 and 128 pixel patches respectively. The top plot corresponds to patch 1 and the bottom one to a patch containing an average value of the projected galaxy density (patch 201). In the three figures we observe that the correlation coefficients found for the patch 1 are consistent which what we expect for tSZ emission with $`A`$ values of $`112\pm 22`$ $`\mu `$K, $`77\pm 16`$ $`\mu `$K and $`54\pm 13`$ $`\mu `$K respectively. Notice that there is good agreement between the model and the data with reduced $`\chi ^2`$ values of $`3/4`$, $`3.4/4`$ and $`2/4`$ for patches containing 32, 64 and 128 pixels respectively. For the null hypothesis ($`A=0`$ $`\mu `$K), the reduced $`\chi ^2`$ values are $`28/5`$, $`26/5`$ and $`20/5`$. We interpret these results as a spectral indication of the measurement of a tSZ signal in patch 1. However, for all other patches we find $`A`$ compatible with zero, showing that the data are in good agreement with the null hypothesis. For example, for patch 201 the $`A`$ values are $`5\pm 19`$ $`\mu `$K, $`4\pm 13`$ $`\mu `$K and $`10\pm 10`$ $`\mu `$K with reduced $`\chi ^2`$ values of $`2.2/4`$, $`1.2/4`$ and $`1.3/4`$ respectively. We consider these results compatible with a no detection of tSZ signal in patch 201. The uncertainties for the total amplitude of the tSZ signal, $`A`$, presented above do not account for systematic errors on the Archeops data. These are dominated by residuals from atmospheric and Galactic dust emissions which in a first order approximation increase as $`\nu ^2`$ in the Archeops frequency range. To account for those contributions we have recomputed the total amplitude for the tSZ signal adding an extra term in the fitted function as follows: $$\alpha (\nu )=A\times f_{tSZ}(\nu )+A_{sys}\times \nu ^2+n(\nu ),$$ (4) where $`A_{sys}`$ is computed only from the Archeops data. For patch 1 we obtain $`A`$ values of $`110\pm 22`$ $`\mu `$K, $`74\pm 16`$ $`\mu `$K and $`53\pm 12`$ $`\mu `$K, and $`A_{sys}`$ values of $`95\pm 102`$ $`\mu `$K, $`104\pm 71`$ $`\mu `$K and $`46\pm 51`$ $`\mu `$K for patches of 32, 64 and 128 pixels respectively. From these results we conclude that the $`A`$ coefficients are not significantly affected at 1-$`\sigma `$ level by systematic effects present in the Archeops data. ### 5.2 Determination of the mean and integrated Compton parameters It is interesting to check whether pixels in patch 1 correspond to potential sources of tSZ or not. For the 32-pixel-size patches, patch 1 includes COMA, A0576, A0671, A0952, A1795, A2061, A2065, A2244, A2245, A2249, A2255. Since these are massive and relatively nearby galaxy clusters, it is reasonable to expect some signature of the tSZ effect. Using the quoted values for the $`A`$ parameter we infer the following estimate for the average Compton parameter in all those sources: $`y=(0.41\pm 0.08)\times 10^4`$. In the case of 64-pixel-size patches we have roughly the same clusters and we expect therefore the signal to be diluted, $`y=(0.28\pm 0.06)\times 10^4`$. Finally for 128 pixel-size patches, patch 1 includes the following 27 clusters : COMA, A0077, A0104, A0272, A0376, A0407, A0576, A0671, A0952, A1035, A1185, A1235, A1377, A1767, A1795, A1800, A2034, A2061, A2065, A2069, A2142, A2151, A2199, A2244, A2245, A2249, A2255. We deduce for them $`y=(0.20\pm 0.05)\times 10^4`$. Since we have scaled the $`\alpha `$’s to the beam-size of the 143 GHz channel of Archeops, these quoted values of $`y`$ are associated to a (linear) angular scale on the sky of $``$ 13 arcmins. For the above results we have assumed the CMB temperature to be $`T_{CMB}=2.725`$K Mather et al. 1999 . We now try to relate the observed average tSZ decrement $`y_{obs}`$ to the high end of the SZ number counts. We make 2 hypotheses: 1) that the Archeops beam encloses most of the integrated tSZ effect in clusters and 2) that the 2MASS survey is a perfect tracer of the tSZ effect. We will come back to these hypotheses at the end. The correlation analysis, that is presented in the previous section, can be recast in stating that the $`N_{pix}`$ brightest clusters of galaxies have an average integrated $`Y=𝑑\mathrm{\Omega }y`$ parameter equal to $`Y=y_{obs}\mathrm{\Omega }_{beam}`$ where $`\mathrm{\Omega }_{beam}=2\pi (\mathrm{FWHM}/\sqrt{8\mathrm{log}2})^2`$ is the Archeops beam solid angle if $`\mathrm{FWHM}=13\mathrm{arcmin}`$. The integrated Compton parameter can be directly related to the SZ flux number counts which has been the issue of a great number of studies (e.g. Aghanim et al. (1997); Bartlett et al. 1994 ; Barbosa et al. 1996 ; Bartelmann et al. 2001 ) from which we select Xue & Wu 2001 and Benson et al. 2002 . Following Xue & Wu 2001 there are three possible models for which the number counts of clusters over the whole sphere can be parametrized in terms of $`Y`$ as follows, $$N(>Y)=N_0(Y/Y_0)^\gamma ,$$ (5) where we fiducially consider $`Y_0=10^2\mathrm{arcmin}^2`$ For Model 1 (M1), deduced from the cosmological $`\mathrm{\Lambda }`$CDM matter power spectrum, $`N_0=29`$ and $`\gamma =1.5`$. For Model 2 (M2), a non-evolving X-ray luminosity function is used to correct the counts and gives a larger $`N_0=635`$ with the same exponent. Finally Model 3 (M3), an evolving X-ray luminosity function is used instead, giving $`N_0=350`$ with a flatter exponent $`\gamma =1.2`$. The conversion from flux in Jansky units to the Compton frequency-independent quantity $`Y`$ is obtained via $$\frac{F_\nu }{Y}=T_{CMB}\frac{B_\nu }{T_{CMB}}f_{tSZ}(\nu ).$$ (6) For example, the fiducial value $`Y_0=10^2arcmin^2`$ is equivalent to 0.75 and 0.91 Jy at 90 and 143 GHz (resp.). From the above formulas, we deduce the lower limit $`Y_{min}`$ in the number counts for the first $`N_{cl}=32,64,128`$ brightest clusters as $$Y_{min}=Y_0\left(\frac{N_{cl}}{N_0f_{sky}}\right)^{1/\gamma }$$ (7) where $`f_{sky}0.20`$ is the effective sky fraction used in the Archeops–WMAP tSZ cross analysis. The average $`Y`$ value of these clusters is then $$Y=Y_0\frac{\gamma }{\gamma 1}\left(\frac{N_{cl}}{N_0f_{sky}}\right)^{1/\gamma }$$ (8) We deduce the average Compton parameter for those clusters $`y`$ as $$y=\frac{Y}{2\pi (\mathrm{FWHM}/\sqrt{8\mathrm{log}2})^2}$$ (9) For the previous three models, expected values for $`y`$ are between $`0.5\times 10^4`$ and $`6\times 10^4`$. The value expected from basic principles (M1), $`y=0.49\times 10^4`$, is quite close to the observed value $`y_{obs}=0.41\pm 0.08\times 10^4`$. The observed value should however be corrected up-wards. Non–linearities in the relation between the 2–MASS density field and the $`Y`$ parameter introduce an efficiency which is difficult to estimate, although we note that as discussed above most of the brightest pixels in the density map constructed from the 2–MASS survey are associated to well-known clusters of galaxies. If we considered that only identified clusters (11) are present in the first patch of 32 pixels, the above expected values should be multiplied by a factor of two. Extension of the bright clusters beyond the fiducial beam of 13 arcmin (like Coma) may produce a differential effect with the observing beam at different frequencies, as well as a bias in the number counts. We also note that the dependence of $`y`$ with the number of pixels taken in the analysis (32, 64 or 128) follows quite well the Model 1 prediction as shown in Table 1. We have also cross-checked the Xue & Wu 2001 modeling with that from Benson et al. 2002 . This alternative model produces an intermediate value between models M1 and M2 and is marginally compatible with the observations. ## 6 Conclusion In this paper we have presented a joint cross-correlation analysis of the Archeops and WMAP data sets with a template of galaxy density constructed from the 2MASS XSC galaxy catalog. We have first divided the 2MASS sky density map in patches of equal number of pixels and sorted these patches in terms of decreasing projected galaxy density. For each of these patches we have performed a cross-correlation analysis with the Archeops data at 143 and 217 GHz and with the WMAP data for the Q, V and W bands. For patches containing the densest 32, 64 and 128 pixels (patch 1), the correlation test pointed to a prominent temperature decrement in WMAP’s Q, V and W bands and in the 143 GHz band of Archeops, but not at 217 GHz, as would be expected for tSZ-induced temperature fluctuations. All the other patches failed to show a similar behavior. To corroborate these results, for each of the patches we have compared the cross-correlation coefficients to a model of the tSZ frequency pattern, in which we fix the spectral behavior and fit for a global amplitude parameter $`A`$. For the first patch and the three sizes considered (32, 64 and 128 pixels), we obtain non zero $`A`$ values at more than 4.5-$`\sigma `$ level with good agreement between the model and the data, and negligible contribution from systematic residuals in the Archeops data. For all other pixels having smaller projected galaxy density, we fail to detect any signature of tSZ effect. From these results, we conclude that there is clear indication of tSZ effect for patch 1 in the Archeops and WMAP data sets. This is not surprising, since patch 1 contains pixels centered in massive and relatively nearby galaxy clusters. Assuming the signal observed is tSZ, we infer, for the 32-pixels case, an average value for the comptonization parameter of $`y=(0.41\pm 0.08)\times 10^4`$ in all those clusters at an angular scale of $``$ 13 arcmins. This value is compatible at 1-$`\sigma `$ level with the expectations, $`y=0.49\times 10^4`$ (cf. Table 1), from a model of the cluster flux number counts based on the standard $`\mathrm{\Lambda }`$-CDM model, M1, assuming the measured $`y`$ is due to the contribution from the 32 brightest clusters. For 64 and 128-pixeles patches the tSZ signal is diluted, probably due to the contribution of relatively not so massive clusters. Note that the dilution observed is also compatible with the one expected from the model M1 (cf. Table 1). ###### Acknowledgements. We acknowledge the Archeops collaboration for the use of the proprietary Archeops data and related software as well as for fruitful comments and careful reading of this manuscript. We acknowledge very useful comments by J.A.Rubiño–Martín. C.H.M. acknowledges the financial support from the European Community through the Human Potential Programme under contract HPRN-CT-2002-00124 (CMBNET). Some of the results in this paper have been derived using the HEALPix package, Górski et al. 1999 . We acknowledge the use of the Legacy Archive for Microwave Background Data Analysis (LAMBDA, http://lambda.gsfc.nasa.gov). Support for LAMBDA is provided by the NASA Office of Space Science. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation.
warning/0506/math0506432.html	ar5iv	text	# The geometry of continued fractions and the topology of surface singularities ## 1. Introduction Continued fraction expansions appear naturally when one resolves germs of plane curves by sequences of plane blowing-ups, or Hirzebruch-Jung (that is, cyclic quotient) surface singularities by toric modifications. They also appear when one passes from the natural plumbing decomposition of the abstract boundary of a normal surface singularity to its minimal JSJ decomposition. In this case it is very important to keep track of natural orientations. In general, as was shown by Neumann , if one changes the orientation of the boundary, the resulting 3-manifold is no more orientation-preserving diffeomorphic to the boundary of an isolated surface singularity. The only exceptions are Hirzebruch-Jung singularities and cusp-singularities. This last class of singularities got its name from its appearance in Hirzebruch’s work as germs at the compactified cusps of Hilbert modular surfaces. For both classes of singularities, one gets an involution on the set of analytical isomorphism types of the singularities in the class, by changing the orientation of the boundary. From the viewpoint of computations, Hirzebruch saw that both types of singularities have structures which can be encoded in continued fraction expansions of positive integers, and that the previous involution manifests itself in a duality between such expansions. In the computations with continued fractions alluded to before, there appear in fact two kinds of continued fraction expansions. Some are constructed using only additions - we call them in the sequel Euclidean continued fractions \- and the others using only subtractions - we call them Hirzebruch-Jung continued fractions. There is a simple formula, also attributed to Hirzebruch, which allows to pass from one type of continued fraction expansion of a number to the other one. Both types of expansions have geometric interpretations in terms of polygonal lines $`P(\sigma )`$. If $`(L,\sigma )`$ is a pair consisting of a 2-dimensional lattice $`L`$ and a strictly convex cone $`\sigma `$ in the associated real vector space, $`P(\sigma )`$ denotes the boundary of the convex hull of the set of lattice points situated inside $`\sigma `$ and different from the origin. For Euclidean continued fractions this interpretation is attributed to Klein , while for the Hirzebruch-Jung ones it is attributed to Cohn . It is natural to try to understand how both geometric interpretations fit together. By superimposing the corresponding drawings, we were led to consider two supplementary cones in a real plane, in the presence of a lattice. By supplementary cones we mean two closed strictly convex cones which have a common edge and whose union is a half-plane. Playing with examples made us understand that the algebraic duality between continued fractions alluded to before has as geometric counterpart a duality between two supplementary cones in the plane with respect to a lattice. This duality is easiest to express in the case where the edges of the cones are irrational: Suppose that the edges of the supplementary cones $`\sigma `$ and $`\sigma ^{}`$ are irrational. Then the edges of each polygonal line $`P(\sigma )`$ and $`P(\sigma ^{})`$ correspond bijectively in a natural way to the vertices of the other one. When at least one of the edges is rational, the correspondence is slightly more complicated (there is a defect of bijectivity near the intersection points of the polygonal lines with the edges of the cones), as explained in Proposition 5.3. In this duality, points correspond to lines and conversely, as in the classical polarity relation between points and lines with respect to a conic. But the duality relation described in this paper is more elementary, in the sense that it uses only parallel transport in the plane. For this reason it can be explained very simply by drawing on a piece of cross-ruled paper. The duality between supplementary cones gives a simple way to think about the relation between the pair $`(L,\sigma )`$ and its dual pair $`(\stackrel{ˇ}{L},\stackrel{ˇ}{\sigma })`$ and in particular about the relations between various invariants of toric surfaces (see section 6). Indeed (see Proposition 5.11): The supplementary cone of $`\sigma `$ is canonically isomorphic over the integers with the dual cone $`\stackrel{ˇ}{\sigma }`$, once an orientation of $`L`$ is fixed. As stated at the beginning of the introduction, computations with continued fractions appear also when one passes from the canonical plumbing structure on the boundary of a normal surface singularity to its minimal JSJ structure. Using this, Neumann showed that the topological type of the minimal good resolution of the germ is determined by the topological type of the link. In fact all continued fractions appearing in Neumann’s work are the algebraic counterpart of pairs $`(L,\sigma )`$ canonically determined by the topology of the boundary. Using this remark, we prove the stronger statement (see Theorem 9.7): The plumbing structure on the boundary of a normal surface singularity associated to the minimal normal crossings resolution is determined up to isotopy by the oriented ambient manifold. In particular, it is invariant up to isotopy under the group of orientation-preserving self-diffeomorphisms of the boundary. In order to prove this theorem we have to treat separately the boundaries of Hirzebruch-Jung and cusp singularities. In both cases, we show that the oriented boundary determines naturally a pair $`(L,\sigma )`$ as before. If one changes the orientation of the boundary, one gets a supplementary cone. In this way, the involution defined before on both sets of singularities is a manifestation of the geometric duality between supplementary cones (see Propositions 9.3 and 9.6). For us, the moral of the story we tell in this paper is the following one: If one meets computations with either Euclidean or Hirzebruch-Jung continued fractions in a geometrical problem, it means that somewhere behind is present a natural 2-dimensional lattice $`L`$ and a couple of lines in the associated real vector space. One has first to choose one of the two pairs of opposite cones determined by the four lines and secondly an ordering of the edges of those cones. These choices may be dictated by choices of orientations of the manifolds which led to the construction of the lattice and the cones. So, in order to think geometrically at the computations with continued fractions, recognize the lattice, the lines and the orientation choices. Let us outline now the content of the paper. Someone who is interested only in the algebraic relations between the Euclidean and the Hirzebruch-Jung continued fraction expansions of a number can consult only section 2. If one is also interested in their geometric interpretation, one can read sections 3 and 4. In section 5 we prove geometrically the relations between the two kinds of continued fractions using the duality between supplementary cones described before. We introduce also a new kind of graphical representation which we call the zigzag diagram, allowing to visualize at the same time the algebra and the geometry of the continued fraction expansions of a number. In section 6 we give applications of zigzag diagrams to the algebraic description of special curve and surface singularities, defined using toric geometry. Sections 8 and 9 are dedicated to the study of topological aspects of the links of normal surface singularities, after having recalled in section 7 general facts about Seifert, graph, plumbing and JSJ-structures on 3-manifolds. We think that the new results of the paper are Proposition 5.3, Theorem 9.1 and Theorem 9.7, as well as the very easy Proposition 5.11, which is nevertheless essential in order to understand the relation between dual cones in terms of parallelism, using Proposition 5.3. We wrote this paper having in mind as a potential reader a graduate student who wants to be initiated either to the algebra of surface singularities or to their topology. That is why we tried to communicate basic intuitions, often referring to the references for complete proofs. Acknowledgments. We are very grateful to Friedrich Hirzebruch for the historical comments he sent us, as well as to Paolo Lisca, Andras Némethi, Bernard Teissier and the anonymous referee for their pertinent remarks and suggestions. ## 2. Algebraic comparison of Euclidean and Hirzebruch-Jung continued fractions ###### Definition 2.1. If $`x_1,\mathrm{},x_n`$ are variables, we consider two kinds of continued fractions associated to them: $$[x_1,\mathrm{},x_n]^+:=x_1+\frac{1}{x_2+{\displaystyle \frac{1}{\mathrm{}+{\displaystyle \frac{1}{x_n}}}}}$$ $$[x_1,\mathrm{},x_n]^{}:=x_1\frac{1}{x_2{\displaystyle \frac{1}{\mathrm{}{\displaystyle \frac{1}{x_n}}}}}$$ We call $`[x_1,\mathrm{},x_n]^+`$ a Euclidean continued fraction (abbreviated E-continued fraction) and $`[x_1,\mathrm{},x_n]^{}`$ a Hirzebruch-Jung continued fraction (abbreviated HJ-continued fraction). The first name is motivated by the fact that E-continued fractions are tightly related to the Euclidean algorithm: if one applies this algorithm to a couple of positive integers $`(a,b)`$ and the successive quotients are $`q_1,\mathrm{},q_n`$, then $`a/b=[q_1,\mathrm{},q_n]^+`$. See Hardy & Wright , Davenport for an introduction to their arithmetics and Fowler for the relation with the Greek theories of proportions. An extended bibliography on their applications can be found in Brezinski and Shallit . The second name is motivated by the fact that HJ-continued fractions appear naturally in the Hirzebruch-Jung method of resolution of singularities, originating in Jung and Hirzebruch , as explained after Definition 6.4 below. Define two sequences $`(Z^\pm (x_1,\mathrm{},x_n))_{n1}`$ of polynomials with integer coefficients, by the initial data $$Z^\pm (\mathrm{})=1,Z^\pm (x)=x$$ and the recurrence relations: (1) $$Z^\pm (x_1,\mathrm{},x_n)=x_1Z^\pm (x_2,\mathrm{},x_n)\pm Z^\pm (x_3,\mathrm{},x_n),n2.$$ Then one proves immediately by induction on $`n`$ the following equality of rational fractions: (2) $$[x_1,\mathrm{},x_n]^\pm =\frac{Z^\pm (x_1,\mathrm{},x_n)}{Z^\pm (x_2,\mathrm{},x_n)},n1.$$ Also by induction on $`n`$, one proves the following twin of relation (1): (3) $$Z^\pm (x_1,\mathrm{},x_n)=Z^\pm (x_1,\mathrm{},x_{n1})x_n\pm Z^\pm (x_1,\mathrm{},x_{n2}),n2.$$ which, combined with (1), proves the following symmetry property: (4) $$Z^\pm (x_1,\mathrm{},x_n)=Z^\pm (x_n,\mathrm{},x_1),n1.$$ If $`(y_1,\mathrm{},y_k)`$ is a finite sequence of numbers or variables and $`m𝐍\{+\mathrm{}\}`$, we denote by $$(y_1,\mathrm{},y_k)^m$$ the sequence obtained by repeating $`m`$ times the sequence $`(y_1,\mathrm{},y_k)`$. By convention, when $`m=0`$, the result is the empty sequence. Each number $`\lambda 𝐑`$ can be expanded as (possibly infinite) Euclidean and Hirzebruch-Jung continued fractions: $$\lambda =[a_1,a_2,\mathrm{}]^+=[\alpha _1,\alpha _2,\mathrm{}]^{}$$ with the conditions: (5) $$a_1𝐙,a_n𝐍\{0\},n1$$ (6) $$\alpha _1𝐙,\alpha _n𝐍\{0,1\},n1$$ Of course, we consider only indices $`n`$ effectively present. For an infinite number of terms, these conditions ensure the existence of the limits $`[a_1,a_2,\mathrm{}]^+:=\underset{n+\mathrm{}}{lim}[a_1,\mathrm{},a_n]^+`$ and $`[\alpha _1,\alpha _2,\mathrm{}]^{}:=\underset{n+\mathrm{}}{lim}[\alpha _1,\mathrm{},\alpha _n]^{}`$. Any sequence $`(a_n)_{n1}`$ which verifies the restrictions (5) can appear and the only ambiguity in the expansion of a number as a E-continued fraction comes from the identity: (7) $$[a_1,\mathrm{},a_n,1]^+=[a_1,\mathrm{},a_{n1},a_n+1]^+$$ We deduce that any real number $`\lambda 1`$ admits a unique expansion as a E-continued fraction such that condition (5) is satisfied and in the case that the sequence $`(a_n)_n`$ is finite, its last term is different from $`1`$. When we speak in the sequel about the E-continued fraction expansion of a number $`\lambda 1`$, it will be about this one. By analogy with the vocabulary of the Euclidean algorithm, we say that the numbers $`(a_n)_{n1}`$ are the E-partial quotients of $`\lambda `$. Similarly, any sequence $`(\alpha _n)_{n1}`$ which verifies the restrictions (6) can appear and the only ambiguity in the expansion of a number as a HJ-continued fraction comes from the identity: (8) $$[\alpha _1,\mathrm{},\alpha _n,(2)^{\mathrm{}}]^{}=[\alpha _1,\mathrm{},\alpha _{n1},\alpha _n1]^{}$$ We see that any real number $`\lambda `$ admits a unique expansion as a HJ-continued fraction such that condition (6) is satisfied and the sequence $`(\alpha _n)_n`$ is not infinite and ultimately constant equal to $`2`$. When we speak in the sequel about the HJ-continued fraction expansion of a number $`\lambda `$, it will be about this one. We call the numbers $`(\alpha _n)_{n1}`$ the HJ-partial quotients of $`\lambda `$. The following lemma (see Hirzebruch \[36, page 257\]) can be easily proved by induction on the integer $`b1`$. ###### Lemma 2.2. If $`a𝐙,b𝐍\{0\}`$ and $`x`$ is a variable, then: $$[a,b,x]^+=[a+1,(2)^{b1},x+1]^{}$$ Using this lemma one sees how to pass from the E-continued fraction expansion of a real number $`\lambda `$ to its HJ-continued fraction expansion: ###### Proposition 2.3. If $`(a_n)_{n1}`$ is a (finite or infinite) sequence of positive integers, then: $$[a_1,\mathrm{},a_{2n}]^+=[a_1+1,(2)^{a_21},a_3+2,(2)^{a_41},\mathrm{},(2)^{a_{2n}1}]^{}$$ $$[a_1,\mathrm{},a_{2n+1}]^+=[a_1+1,(2)^{a_21},a_3+2,(2)^{a_41},\mathrm{},(2)^{a_{2n}1},a_{2n+1}+1]^{}$$ $$[a_1,a_2,a_3,a_4,\mathrm{}]^+=[a_1+1,(2)^{a_21},a_3+2,(2)^{a_41},a_5+2,(2)^{a_61},\mathrm{}]^{}$$ (recall that, by convention, $`(2)^0`$ denotes the empty sequence). ###### Example 2.4. $`\frac{11}{7}=[(1)^3,3]^+=[2,3,(2)^2]^{}`$. Notice that this procedure can be inverted. In particular, an immediate consequence of the previous proposition is that a number has bounded E-partial quotients if and only if it has bounded HJ-partial quotients. Similarly, it has ultimately periodic E-continued fraction (which happens if and only if it is a quadratic number, see Davenport ) if and only if it has ultimately periodic HJ-continued fraction. In this case, Proposition 2.3 explains how to pass from its E-period to its HJ-period. The continued fraction expansions of two numbers which differ by an integer are related in an evident and simple way. For this reason, from now on we restrict our attention to real numbers $`\lambda >1`$. The map (9) $$\lambda \frac{\lambda }{\lambda 1}$$ is an involution of the interval $`(1,+\mathrm{})`$ on itself. The E-continued fraction expansions of the numbers in the same orbit of this involution are related in a very simple way: ###### Lemma 2.5. If $`\lambda (1,+\mathrm{})`$ and $`\lambda =[a_1,a_2,\mathrm{}]^+`$ is its expansion as a (finite or infinite) continued fraction, then: $$\frac{\lambda }{\lambda 1}=\{\begin{array}{cc}[1+a_2,a_3,a_4,\mathrm{}]^+,\hfill & \text{ if }a_1=1,\hfill \\ [1,a_11,a_2,a_3,\mathrm{}]^+,\hfill & \text{ if }a_12\hfill \end{array}$$ The proof is immediate, once one notices that $`\frac{\lambda }{\lambda 1}=[1,\lambda 1]^+`$. Notice also that the involutivity of the map (9) shows that the first equality in the previous lemma is equivalent to the second one. ###### Example 2.6. If $`\lambda =\frac{11}{7}=[(1)^3,3]^+`$, then $`\frac{11}{4}=\frac{\lambda }{\lambda 1}=[2,1,3]^+`$. By combining Proposition 2.3 and Lemma 2.5, we get the following relation between the HJ-continued fraction expansions of the numbers in the same orbit of the involution (9): ###### Proposition 2.7. If $`\lambda 𝐑`$ is greater than $`1`$ and $$\lambda =[(2)^{m_1},n_1+3,(2)^{m_2},n_2+3,\mathrm{}]^{}$$ is its expression as a (finite or infinite) continued fraction, with $`m_i,n_i𝐍,i1`$, then: $$\frac{\lambda }{\lambda 1}=[m_1+2,(2)^{n_1},m_2+3,(2)^{n_2},m_3+3,\mathrm{}]^{}$$ For $`\lambda `$ rational, this was proved in a different way by Neumann \[56, Lemma 7.2\]. It reads then: $$\begin{array}{cc}\lambda =[(2)^{m_1},n_1+3,(2)^{m_2},\mathrm{},n_s+3,(2)^{m_{s+1}}]^{}\hfill & \\ \frac{\lambda }{\lambda 1}=[m_1+2,(2)^{n_1},m_2+3,\mathrm{},(2)^{n_s},m_{s+1}+2]^{}\hfill & \end{array}$$ The important point here is that even a value $`m_{s+1}=0`$ contributes to the number of partial quotients in the HJ-continued fraction expansion of $`\frac{\lambda }{\lambda 1}`$. The next proposition is equivalent to the previous one, as an easy inspection shows. Its advantage is that it gives a graphical way to pass from the HJ-continued fraction expansion of a number $`\lambda >1`$ to the analogous expansion of $`\frac{\lambda }{\lambda 1}>1`$. ###### Proposition 2.8. Consider a number $`\lambda 𝐑`$ greater than $`1`$ and let $$\lambda =[\alpha _1,\alpha _2,\mathrm{}]^{},\frac{\lambda }{\lambda 1}=[\beta _1,\beta _2,\mathrm{}]^{}$$ be the expressions of $`\lambda `$ and $`\frac{\lambda }{\lambda 1}`$ as (finite or infinite) HJ-continued fractions. Construct a diagram made of points organized in lines and columns in the following way: $``$ its lines are numbered by the positive integers; $``$ the line numbered $`k1`$ contains $`\alpha _k1`$ points; $``$ the first point in the line numbered $`k+1`$ is placed under the last point of the line numbered $`k`$. Then the $`k`$-th column contains $`\beta _k1`$ points. This graphical construction seems to have been first noticed by Riemenschneider in when $`\lambda 𝐐_+`$. Nowadays one usually speaks about Riemenschneider’s point diagram or staircase diagram. ###### Example 2.9. If $`\lambda =\frac{11}{7}=[2,3,(2)^2]^{}`$, the associated point diagram is: $$\begin{array}{cc}& \\ & \\ & \\ & \end{array}$$ One deduces from it that $`\frac{\lambda }{\lambda 1}=[3,4]^{}`$. ## 3. Klein’s geometric interpretation of Euclidean continued fractions We let Klein himself speak about his interpretation, in order to emphasize his poetical style: > Let us now enliven these considerations with geometric pictures. Confining our attention to positive numbers, let us mark all those points in the positive quadrant of the $`xy`$ plane which have integral coordinates, forming thus a so-called point lattice. Let us examine this lattice, I am tempted to say this “firmament” of points, with our point of view at the origin. \[…\] Looking from $`0`$, then, one sees points of the lattice in all rational directions and only in such directions. The field of view is everywhere “densely” but not completely and continuously filled with “stars”. One might be inclined to compare this view with that of the milky way. With the exception of $`0`$ itself there is not a single integral point lying upon an irrational ray $`\frac{x}{y}=\omega `$, where $`\omega `$ is irrational, which is very remarkable. If we recall Dedekind’s definition of irrational number, it becomes obvious that such a ray makes a cut in the field of integral points by separating the points into two point sets, one lying to the right of the ray and one to the left. If we inquire how these point sets converge toward our ray $`x/y=\omega `$, we shall find a very simple relation to the continued fraction for $`\omega `$. By marking each point $`(x=p_\nu ,y=q_\nu )`$, corresponding to the convergent $`p_\nu /q_\nu `$, we see that the rays to these points approximate to the ray $`x/y=\omega `$ better and better, alternately from the left and from the right, just as the numbers $`p_\nu /q_\nu `$ approximate to the number $`\omega `$. Moreover, if one makes use of the known number-theoretic properties of $`p_\nu ,q_\nu `$, one finds the following theorem: Imagine pegs or needles affixed at all the integral points, and wrap a tightly drawn string about the sets of pegs to the right and to the left of the $`\omega `$-ray, then the vertices of the two convex string-polygons which bound our two point sets will be precisely the points $`(p_\nu ,q_\nu )`$ whose coordinates are the numerators and denominators of the successive convergents to $`\omega `$, the left polygon having the even convergents, the right one the odd. This gives a new and, one may well say, an extremely graphical definition of a continued fraction. In the original article , one finds moreover the following interpretation of the E-partial quotients: > Each edge of the polygons \[…\] may contain integral points. The number of parts in which the edge is decomposed by such points is exactly equal to a partial quotient. Before Klein, Smith expressed a related idea in : > If with a pair of rectangular axes in a plane we construct a system of unit points (i.e. a system of points of which the coordinates are integral numbers), and draw the line $`y=\theta x`$, we learn from that theorem that if $`(x,y)`$ be a unit point lying nearer to that line than any other unit point having a less abscissa (or, which comes to the same thing, lying at a less distance from the origin), $`\frac{y}{x}`$ is a convergent to $`\theta `$; and, vice versa, if $`\frac{y}{x}`$ is a convergent, $`(x,y)`$ is one of the ‘nearest points’. Thus the ‘nearest points’ lie alternately on opposite sides of the line, and the double area of the triangle, formed by the origin and any two consecutive ‘nearest points’, is unity. Proofs of the preceding properties can be found in Stark . Here we only sketch the reason of Klein’s interpretation. For explanations about our vocabulary, read next section. Let $`\lambda >1`$ be a real number. In the first quadrant $`\sigma _0`$, consider the half-line $`L_\lambda `$ of slope $`\lambda `$ (see Figure 1). It is defined by the equation $`y=\lambda x`$, which shows that $`\lambda =\omega ^1=\theta `$, where $`\omega `$ is Klein’s notation and $`\theta `$ is Smith’s. It subdivides the quadrant $`\sigma _0`$ into two closed cones with vertex the origin, $`\sigma _x(\lambda )`$ adjacent to the axis of the variable $`x`$ and $`\sigma _y(\lambda )`$ adjacent to the axis of the variable $`y`$. ###### Lemma 3.1. The segment which joins the lattice points of coordinates $`(1,0)`$ and $`(1,a_1)`$ is a compact edge of the convex hull of the set of lattice points different from the origin contained in the cone $`\sigma _x(\lambda )`$, where $`\lambda =[a_1,a_2,\mathrm{}]^+`$ is the E-continued fraction expansion of $`\lambda `$. Proof: Indeed, the half-line starting from $`(1,0)`$ and directed towards $`(1,a_1)`$ cuts the half-line $`L_\lambda `$ inside the segment $`[(1,[\lambda ]),(1,[\lambda ]+1))`$, where $`[\lambda ]`$ is the integral part of $`\lambda `$. But $`[\lambda ]=a_1`$, which finishes the proof. $`\mathrm{}`$ Replace now the initial basis of the lattice by $`(0,1),(1,a_1)`$. With respect to this new basis, the slope of the half-line $`L_\lambda `$ is $`(\lambda a_1)^1=[a_2,a_3,\mathrm{}]^+`$. This allows one to prove Klein’s interpretation by induction. If one considers all lattice points on the compact edges of the boundaries of the two previous convex hulls instead of only the vertices, and then one looks at the slopes of the lines which join them to the origin, one obtains the so-called slow approximating sequence of $`\lambda `$. This kind of sequence appears naturally when one desingularizes germs of complex analytic plane curves by successively blowing up points (see Enriques & Chisini , Michel & Weber and Lê, Michel & Weber ). We leave as an exercise for the interested reader to interpret this geometrically (first, read Section 6.3). As explained by Klein himself in , his interpretation suggests to generalize the notion of continued fraction to higher dimensions by taking the boundaries of convex hulls of lattice points situated inside convex cones. For references about recent research in this area, see Arnold and Moussafir . ## 4. Cohn’s geometric interpretation of Hirzebruch-Jung continued fractions A geometric interpretation of HJ-continued fractions analogous to Klein’s interpretation of Euclidean ones was given by Cohn (see the comment on his work in Hirzebruch \[36, 2.3\]). It seems to have soon become folklore among people doing toric geometry (see section 6). Before describing this interpretation, let us introduce some vocabulary in order to speak with more precision about convex hulls of lattice points in the plane. Let $`L`$ be a lattice of rank $`2`$, that is, a free abelian group of rank $`2`$. It embeds canonically into the associated real vector space $`L_𝐑=L_𝐙𝐑`$. When we picture the elements of $`L`$ as points in the affine plane $`L_𝐑`$, we call them the integral points of the plane. When $`A`$ and $`B`$ are points of the affine plane $`L_𝐑`$, we denote by $`AB`$ the element of the vector space $`L_𝐑`$ which translates $`A`$ into $`B`$, by $`[AB]`$ the closed segment in $`L_𝐑`$ of extremities $`A,B`$ and by $`[AB`$ the closed half-line having $`A`$ as an extremity and directed towards $`B`$. If $`(u,v)`$ is an ordered basis of $`L_𝐑`$ and $`l`$ is a line of $`L_𝐑`$, its slope is the quotient $`\beta /\alpha 𝐑\{\mathrm{}\}`$, where $`\alpha u+\beta v`$ generates $`l`$. ###### Definition 4.1. A (closed convex) triangle $`ABC`$ in $`L_𝐑`$ is called elementary if its vertices are integral and they are the only intersections of the triangle with the lattice $`L`$. If the triangle $`ABC`$ is elementary, then each pair of vectors $`(AB,AC)`$, $`(BC,BA)`$, $`(CA,CB)`$ is a basis of the lattice $`L`$. Conversely, if one of these pairs is a basis of the lattice, then the triangle is elementary. We call a line or a half-line in $`L_𝐑`$ rational if it contains at least two integral points. If so, then it contains an infinity of them. If $`A`$ and $`B`$ are two integral points, the integral length $`l_𝐙[AB]`$ of the segment $`[AB]`$ is the number of subsegments in which it is divided by the integral points it contains. A vector $`OA`$ of $`L`$ is called primitive if $`l_𝐙[OA]=1`$. Let $`\sigma `$ be a closed strictly convex 2-dimensional cone in the plane $`L_𝐑`$, that is, the convex “angle” (in the language of plane elementary geometry) delimited by two non-opposing half-lines originating from $`0`$. These half-lines are called the edges of $`\sigma `$. The cone $`\sigma `$ is called rational if its edges are rational. A cone is called regular if its edges contain points $`A,B`$ such that the triangle $`OAB`$ is elementary. The name is motivated by the fact that the associated toric surface $`𝒵(L,\sigma )`$ is smooth (that is, all its local rings are regular) if and only if $`\sigma `$ is regular (see section 6.1). Let $`K(\sigma )`$ be the convex hull of the set of lattice points situated inside $`\sigma `$, with the exception of the origin, that is: $$K(\sigma ):=\mathrm{Conv}(\sigma (L\{0\})).$$ The closed convex set $`K(\sigma )`$ is unbounded. Denote by $`P(\sigma )`$ its boundary: it is a connected polygonal line. It has two ends (in the topological sense), each one being asymptotic to (or contained inside) an edge of $`\sigma `$ (see Figure 2). An edge of $`\sigma `$ intersects $`P(\sigma )`$ if and only if it is rational. Denote by $`𝒱(\sigma )`$ the set of vertices of $`P(\sigma )`$ and by $`(\sigma )`$ the set of its (closed) edges. For example, in Figure 3 the vertices are the points $`A_0,A_2,A_5`$ and the edges are the segments $`[A_0A_2],[A_2A_5]`$ and two half-lines contained in $`l_{},l_+`$, starting from $`A_0`$, respectively $`A_5`$. Now order arbitrarily the edges of $`\sigma `$. Denote by $`l_{}`$ the first one and by $`l_+`$ the second one. This orients the plane $`L_𝐑`$, by deciding to turn from $`l_{}`$ towards $`l_+`$ inside $`\sigma `$. If we orient $`P(\sigma )`$ from the end which is asymptotic to $`l_{}`$ towards the end which is asymptotic to $`l_+`$, we get induced orientations of its edges. Suppose now that the edge $`l_{}`$ of $`\sigma `$ is rational. Denote then by $`A_{}0`$ the integral point of the half-line $`l_{}`$ which lies nearest to $`0`$, and by $`V_{}A_{}`$ the vertex of $`P(\sigma )`$ which lies nearest to $`A_{}`$. Define in the same way $`A_+`$ and $`V_+`$ whenever $`l_+`$ is rational. Denote by $`(A_n)_{n0}`$ the sequence of integral points on $`P(\sigma )`$, enumerated as they appear when one travels on this polygonal line in the positive direction, starting from $`A_0=A_{}`$. If $`l_+`$ is a rational half-line, then we stop this sequence when we arrive at the point $`A_+`$. If $`l_+`$ is irrational, then this sequence is infinite. Define $`r0`$ such that $`A_{r+1}=A_+`$. So, $`r=+\mathrm{}`$ if and only if $`l_+`$ is irrational. ###### Example 4.2. We consider the lattice $`𝐙^2𝐑^2`$ and the cone $`\sigma `$ with rational edges, generated by the vectors $`(1,0)`$ and $`(4,11)`$ (see Figure 3). The small dots represent integral points in the plane and the bigger ones represent integral points on the polygonal lines $`P(\sigma )`$. In this example we have $`V_+=V_{}=A_2`$. Each triangle $`OA_nA_{n+1}`$ is elementary, by the construction of the convex hull $`K(\sigma )`$, which implies that all the couples $`(OA_n,OA_{n+1})`$ are bases of $`L`$. This shows that for any $`n\{1,\mathrm{},r\}`$, one has a relation of the type: (10) $$OA_{n+1}+OA_{n1}=\alpha _nOA_n$$ with $`\alpha _n𝐙`$, and the convexity of $`K(\sigma )`$ shows that: (11) $$\alpha _n2,n\{1,\mathrm{},r\}$$ Conversely: ###### Proposition 4.3. Suppose that $`(OA_n)_{n0}`$ is a (finite or infinite) sequence of primitive vectors of $`L`$, related by relations of the form (10). Then we have $$OA_n=Z^{}(\alpha _1,\mathrm{},\alpha _{n1})OA_1Z^{}(\alpha _2,\mathrm{},\alpha _{n1})OA_0,n1$$ and the slope of the half-line $`l_+=lim_n\mathrm{}[OA_n)`$ in the base $`(OA_0,OA_1)`$ is equal to $`[\alpha _1,\alpha _2,\mathrm{}]^{}`$. Proof: Recall that the polynomials $`Z^{}`$ were defined by the recursion formula (1). The first assertion can be easily proved by induction, using the relations (10). The second one is a consequence of formula (2), which shows that the slope of the half-line $`[OA_n`$ in the base $`(OA_0,OA_1)`$ is equal to $`[\alpha _1,\mathrm{},\alpha _{n1}]^{}`$. $`\mathrm{}`$ ###### Proposition 4.4. Let $`\sigma `$ be the closure of the convex hull of the union of the half-lines $`([OA_n)_{n0}`$. Then $`\sigma `$ is strictly convex and the points $`\{A_n\}_{n1}`$ are precisely the integral points on the compact edges of the polygonal line $`P(\sigma )`$ if and only if the conditions (11) are satisfied and the sequence $`(\alpha _n)_{n1}`$ is not infinite and ultimately constant equal to $`2`$. Proof: $``$ What remains to be proved about the necessity is that if the sequence $`(\alpha _n)_{n1}`$ is infinite, then it cannot be ultimately constant equal to $`2`$. If this was the case, by relation (8) we would deduce that $`[\alpha _1,\alpha _2,\mathrm{}]^{}`$ is rational, and Proposition 4.3 would imply that $`l_+`$ is rational. Then $`P(\sigma )`$ would contain a finite number of integral points on its compact edges, which would contradict the infinity of the sequence $`(\alpha _n)_{n1}`$. $``$ Let us prove now the sufficiency. As $`\alpha _n2`$, $`n\{1,\mathrm{},r\}`$, we see that the triangles $`(OA_nA_{n+1})_{n0}`$ turn in the same sense. Moreover, Proposition 4.3 shows that $`\sigma `$ is a strictly convex cone. The vertices of the polygonal line $`P=A_0A_1A_2\mathrm{}`$ are precisely those points $`A_n`$ for which $`\alpha _n3`$. As all the triangles $`OA_nA_{n+1}`$ are elementary, we see that the origin $`O`$ is the only integral point of the connected component of $`\sigma P`$ which contains it. Moreover, conditions (11) show that the other component is convex. So, $`PP(\sigma )`$. The proposition is proved. $`\mathrm{}`$ ## 5. Geometric comparison of Euclidean and HJ-continued fractions In section 5.1 we relate the two preceding interpretations, by describing a duality between two supplementary cones in the plane, an underlying lattice being fixed (see Proposition 5.3). In section 5.2 we introduce a so-called zigzag diagram based on this duality, which makes it very easy to visualize the various relations between continued fractions proved algebraically in section 2. In section 5.3 we give a proof of the isomorphism between the supplementary cone $`(L,\sigma ^{})`$ and the dual cone $`(\stackrel{ˇ}{L},\stackrel{ˇ}{\sigma })`$ of a given cone $`(L,\sigma )`$. ### 5.1. A geometric duality between supplementary cones Suppose again that $`\sigma `$ is any strictly convex cone in $`L_𝐑`$, whose edge $`l_{}`$ is not necessarily rational. Let $`l_{}^{}`$ be the half-line opposite to $`l_{}`$ and $`\sigma ^{}`$ be the closed convex cone bounded by $`l_+`$ and $`l_{}^{}`$. So, $`\sigma `$ and $`\sigma ^{}`$ are supplementary cones: ###### Definition 5.1. Two strictly convex cones in a real plane are called supplementary if they have a common edge and if their union is a half-plane. By analogy with what we did in the previous section for $`\sigma `$, orient the polygonal line $`P(\sigma ^{})`$ from $`l_{}^{}`$ towards $`l_+`$. If $`l_{}`$ is rational, define the point $`A_{}^{}`$ and the sequence $`(A_n^{})_{n0}`$, with $`A_0^{}=A_{}^{}`$. They are the analogs for $`\sigma ^{}`$ of the points $`A_{}`$ and $`(A_n)_{n0}`$ for $`\sigma `$. In particular, $`OA_{}+OA_{}^{}=0`$. ###### Example 5.2. Consider the same cone as in Example 4.2. Then the polygonal lines $`P(\sigma )`$ and $`P(\sigma ^{})`$ are represented in Figure 4 using heavy segments. The basis for our geometric comparison of Euclidean and Hirzebruch-Jung continued fractions is the observation that the polygonal line $`P(\sigma ^{})`$ can be constructed in a very simple way once one knows $`P(\sigma )`$. Namely, starting from the origin, one draws the half-lines parallel to the oriented edges of $`P(\sigma )`$. On each half-line, one considers the integer point which is nearest to the origin. Then the polygonal line which joins those points is the union of the compact edges of $`P(\sigma ^{})`$. Now we describe this with more precision. If $`e(\sigma )`$ is an edge of $`P(\sigma )`$, denote by $`(e)L`$ the integral point such that $`O(e)`$ is a primitive vector of $`L`$ positively parallel to $`e`$ (where $`e`$ is oriented according to the chosen orientation of $`P(\sigma )`$). Then it is an easy exercise to see that $`(e)\sigma ^{}`$ (use the fact that the line containing $`e`$ intersects $`l_{}`$ and $`l_+`$ in interior points). We can define a map: (12) $$\begin{array}{cccc}:\hfill & (\sigma )\hfill & & \sigma ^{}L\hfill \\ & e\hfill & & (e)\hfill \end{array}$$ As the edges of $`P(\sigma )`$ always turn in the same direction, one sees that the map $``$ is injective. ###### Proposition 5.3. The map $``$ respects the orientations and the image of $``$ verifies the double inclusion $$𝒱(\sigma ^{})\mathrm{Im}()P(\sigma ^{})L.$$ The difference $`\mathrm{Im}()𝒱(\sigma ^{})`$ contains at most the points $`[A_{}V_{}]`$, $`[V_+A_+]`$. Such a point is a vertex of $`P(\sigma ^{})`$ if and only if the integral length of the corresponding edge of $`P(\sigma )`$ is $`2`$. In particular, one has the equality $`𝒱(\sigma ^{})=\mathrm{Im}()`$ if and only if $`l_𝐙[A_{}V_{}]2`$ and $`l_𝐙[V_+A_+]2`$, whenever these segments exist. Proof: Denote by $`(V_j)_{jJ}`$ the vertices of $`P(\sigma )`$, enumerated in the positive direction. The indices form a set of consecutive integers, well-defined only up to translations. For any $`jJ`$, denote by $`V_j^{}`$ and $`V_j^+`$ respectively the integral points of $`P(\sigma )`$ which precede and follow $`V_j`$. If $`V_j`$ is an interior point of $`\sigma `$, denote by $`W_jL`$ the point such that $`OW_j=OV_j^{}+OV_j^+`$, and by $`W_j^{}`$ its nearest integral point in the interior of the segment $`[OW_j]`$ (see Figure 5). As $`OV_j^{}V_j`$ and $`OV_jV_j^+`$ are elementary triangles, it implies that both $`(OV_j^{},OV_j)`$ and $`(OV_j,OV_j^+)`$ are bases of $`L`$. So, there exists an integer $`n_j`$ such that (13) $$OV_j^{}+OV_j^+=(n_j+3)OV_j.$$ As $`V_j`$ is a vertex of $`P(\sigma )`$, we see that $`n_j0`$. We deduce that the points $`O,V_j,W_j^{},W_j`$ are aligned in this order, that $`V_jV_j^{}+V_jV_j^+=V_jW_j^{}`$ and that $`l_𝐙[V_jW_j^{}]=n_j+1`$. Let us join each one of the $`n_j`$ interior points of $`[V_jW_j^{}]`$ to $`V_j^{}`$. This gives a decomposition of the triangle $`V_j^{}V_jW_j^{}`$ into $`(n_j+1)`$ triangles. These are necessarily elementary, because the triangle $`OV_j^{}V_j`$ is. Denote $$V_j^{}=[V_{j1}V_j]\mathrm{and}V_{j+1}^{}=[V_jV_{j+1}].$$ By the definition of the map $``$, we see that $`OV_j^{}=V_j^{}V_j`$ and $`OV_{j+1}^{}=V_jV_j^+=V_j^{}W_j^{}.`$ This implies that the triangle $`OV_j^{}V_{j+1}^{}`$ is the translated image by the vector $`V_j^{}O`$ of the triangle $`V_j^{}V_jW_j^{}`$. The preceding arguments show that its only integral points are its vertices and $`n_j`$ other points in the interior of the segment $`[V_j^{}V_{j+1}^{}]`$. Indeed: (14) $$V_j^{}V_{j+1}^{}=V_jW_j^{}=(n_j+1)OV_j$$ Moreover, the triangle $`OV_j^{}V_{j+1}^{}`$is included in the cone $`\sigma ^{}`$ and the couple of vectors $`(OV_j^{},OV_{j+1}^{})`$ has the same orientation as $`(l_{}^{},l_+)`$. This shows that the triangles $`(OV_j^{}V_{j+1}^{})_{jJ}`$ are pairwise disjoint and that their union does not contain integral points in its interior. $``$ If both edges of $`\sigma `$ are irrational, then the closure of the union of the cones $`𝐑_+OV_j^{}+𝐑_+OV_{j+1}^{}`$ is the cone $`\sigma ^{}`$, as the edges $`l_{}`$ and $`l_+`$ are asymptotic to $`P(\sigma )`$. We deduce from relation (14) that the sequence $`(\lambda _j)_{jJ}`$ of slopes of the vectors $`(V_j^{}V_{j+1}^{})_{jJ}`$, expressed in a base $`(u_{},u_+)`$ of $`L_𝐑`$ which verifies $`l_\pm =𝐑_+u_\pm `$ is strictly increasing, and that $`lim_j\mathrm{}\lambda _j=0`$, $`lim_{j+\mathrm{}}\lambda _j=+\mathrm{}`$. This shows that the closure of the connected component of $`\sigma ^{}_{jJ}[V_j^{}V_{j+1}^{}]`$ which does not contain the origin is convex. As a consequence, $$\underset{jJ}{}[V_j^{}V_{j+1}^{}]=P(\sigma ^{}).$$ Moreover, as $`n_j0`$, the strict monotonicity of the sequence $`(\lambda _j)_{jJ}`$ implies that the points $`(V_j^{})_{jJ}`$ are precisely the vertices of $`P(\sigma ^{})`$. The proposition is proved in this case. $``$ Suppose now that $`l_{}`$ is rational. Then choose the index set $`J`$ such that $`V_0=A_{}`$ and $`V_1=V_{}`$. By the construction of the map $``$, the triangle $`OV_0^{}V_1^{}`$ is the translated image of $`V_0OV_0^+`$ by the vector $`V_0O`$ (see Figure 6). In particular, $`V_0^{}V_1^{}=OV_0^+`$. But $`V_1^{}V_2^{}=(n_1+1)OV_1`$ by relation (14), which shows that the vectors $`V_0^{}V_1^{}`$ and $`V_1^{}V_2^{}`$ are proportional if and only if $`V_0^+=V_1`$, which is equivalent to $`l_𝐙[A_{}V_{}]=1`$. Moreover, the property of monotonicity for the slopes of the vectors $`(V_j^{}V_{j+1}^{})_{jJ}`$ is true as before, if one starts from $`j=0`$. $``$ An analogous reasoning is valid for $`l_+`$ if this edge of $`\sigma `$ is rational. By combining all this, the proposition is proved. $`\mathrm{}`$ The previous proposition explains a geometric duality between the supplementary cones $`\sigma ,\sigma ^{}`$ with respect to the lattice $`L`$. We see that, with possible exceptions for the compact edges which intersect the edges of $`\sigma `$ and $`\sigma ^{}`$, the compact edges of $`P(\sigma )`$ correspond to the vertices of $`P(\sigma ^{})`$ interior to $`\sigma ^{}`$ and conversely (by permuting the roles of $`\sigma `$ and $`\sigma ^{}`$), which is a kind of point-line polarity relation. The next corollary shows that the involution (9) studied algebraically in section 2 is closely related to the previous duality. ###### Corollary 5.4. Suppose that $`l_{}`$ is rational and that $`\sigma `$ is not regular. If $`(OA_0^{},U)`$ is a basis of $`L`$ with respect to which the slope of $`l_+`$ is greater than $`1`$, then $`U=OA_1`$. If $`\lambda >1`$ denotes the slope of the half-line $`l_+`$ in the base $`(OA_0^{},OA_1)`$, then $`\frac{\lambda }{\lambda 1}`$ is its slope in the base $`(OA_0,OA_1^{})`$. Proof: We leave the first affirmation to the reader (look at Figure 6). As the triangles $`OA_0A_1`$ and $`OA_0^{}A_1^{}`$ are elementary, we see that $`(OA_0,OA_1)`$ and $`(OA_0^{},OA_1^{})`$ are indeed two bases of the lattice $`L`$. Proposition 5.3 shows that $`OA_0^{}=A_0A_1`$, which allows us to relate the two bases: (15) $$\{\begin{array}{c}OA_0^{}=OA_0\hfill \\ OA_1^{}=OA_1OA_0\hfill \end{array}$$ Let $`vL_𝐑`$ be a vector which generates the half-line $`l_+`$. We want to express it in these two bases. As $`l_+`$ lies between the half-lines $`[OA_0^{}`$ and $`[OA_1`$, we see that: (16) $$v=qOA_0+pOA_1,\text{ with }p,q𝐑_+^{}$$ The equations (15) imply then that: (17) $$v=(pq)OA_0^{}+pOA_1^{}$$ which shows that $`pq>0`$, as $`l_+`$ lies between the half-lines $`[OA_1^{}`$ and $`[OA_0`$. This implies that $`\lambda :=\frac{p}{q}>1`$. We then deduce the corollary from equation (17). $`\mathrm{}`$ The previous corollary shows that the number $`\lambda >1`$ can be canonically attached to the pair $`(L,\sigma )`$, once a rational edge of $`\sigma `$ is chosen as the first edge $`l_{}`$. This motivates the following definition: ###### Definition 5.5. Suppose that $`l_{}`$ is rational and that the cone $`\sigma `$ is not regular. We say that the pair $`(L,\sigma )`$ with the chosen ordering of sides is of type $`\lambda >1`$ if $`\lambda `$ is the slope of the half-line $`l_+`$ in the base $`(OA_0^{},OA_1)`$. Proposition 4.3 shows that, if $`(L,\sigma )`$ is of type $`\lambda >1`$, then $`\lambda =[\alpha _1,\alpha _2,\mathrm{}]^{}`$, where the sequence $`(\alpha _n)_{n1}`$ was defined using relation (10). Suppose now that both edges of $`\sigma `$ are rational. Then one can choose $`p,q𝐍^{}`$ with $`gcd(p,q)=1`$ in relation (16), condition which determines them uniquely. So, $`\lambda =\frac{p}{q}`$. The following proposition describes the type of $`(L,\sigma )`$ after changing the ordering of the sides. ###### Proposition 5.6. If $`(L,\sigma )`$ is of type $`{\displaystyle \frac{p}{q}}`$ with respect to the ordering $`l_{},l_+`$, then it is of type $`{\displaystyle \frac{p}{\overline{q}}}`$ with respect to the ordering $`l_+,l_{}`$, where $`q\overline{q}1(\text{mod }p)`$. Proof: By relation (16), we have $`OA_+=qOA_{}+pOA_1`$. Multiply both sides by $`\overline{q}`$. By the definition of $`\overline{q}`$, there exists $`k𝐍`$ such that $`q\overline{q}=1+kp`$. We deduce that $`OA_{}=\overline{q}OA_++p(\overline{q}OA_1kOA_{})`$. So, $`(OA_+,\overline{q}OA_1kOA_{})`$ is a base of $`L`$ in which the slope of $`l_{}`$ is $`\frac{p}{\overline{q}}>1`$. By the first affirmation of Corollary 5.4, the proposition is proved. $`\mathrm{}`$ By combining the previous proposition with Proposition 4.3, we deduce the following classical fact (see \[4, section III.5\]): ###### Corollary 5.7. If $`{\displaystyle \frac{p}{q}}=[\alpha _1,\alpha _2,\mathrm{},\alpha _r]^{}`$, then $`{\displaystyle \frac{p}{\overline{q}}}=[\alpha _r,\alpha _{r1},\mathrm{},\alpha _1]^{}`$. Another immediate consequence of Corollary 5.4 is: ###### Proposition 5.8. If $`(L,\sigma )`$ is of type $`{\displaystyle \frac{p}{q}}`$ with respect to the ordering $`l_{},l_+`$, then $`(L,\sigma ^{})`$ is of type $`{\displaystyle \frac{p}{pq}}`$ with respect to the ordering $`l_{}^{},l_+`$. The previous proposition describes the relation between the types of two supplementary cones. In section 5.2, we describe more precisely the relation between numerical invariants attached to the edges and the vertices of $`P(\sigma )`$ and $`P(\sigma ^{})`$. ### 5.2. A diagram relating Euclidean and HJ-continued fractions We introduce now a diagram which allows one to “see” the duality between $`P(\sigma )`$ and $`P(\sigma ^{})`$, as well as the relations between the various numerical invariants attached to these polygonal lines. $``$ Suppose first that both $`l_{}`$ and $`l_+`$ are irrational. Consider two consecutive vertices $`V_j,V_{j+1}`$ of $`P(\sigma )`$. Let us attach the weight $`n_j+3`$ to the vertex $`V_j`$, where $`n_j0`$ was defined by relation (13). Introduce also the integer $`m_{j+1}0`$ such that $`l_𝐙[V_jV_{j+1}]=m_{j+1}+1`$. The relation (14) shows that $`l_𝐙[V_j^{}V_{j+1}^{}]=n_j+1`$. By reversing the roles of the polygonal lines $`P(\sigma ^{})`$ and $`P(\sigma )`$, we deduce that the weight of the vertex $`V_{j+1}^{}`$ of $`P(\sigma ^{})`$ is $`m_{j+1}+3`$. We can visualize the relations between the vertices $`V_j,V_{j+1},V_j^{},V_{j+1}^{}`$ as well as the numbers associated to them and to the segments $`[V_jV_{j+1}]`$, $`[V_j^{}V_{j+1}^{}]`$ by using a diagram, in which the heavy lines represent the polygonal lines $`P(\sigma ),P(\sigma ^{})`$, and each vertex $`V_j`$ is joined to $`V_j^{}`$ and $`V_{j+1}^{}`$ (see Figure 7). In this way, the region contained between the two curves representing $`P(\sigma )`$ and $`P(\sigma ^{})`$ is subdivided into triangles. Each edge $`E`$ of $`P(\sigma )`$, $`P(\sigma ^{})`$ is contained in only one of those triangles. Look at its opposite vertex. We say that $`E`$ is the opposite edge of that vertex in the zigzag diagram. We see that the weight of a vertex is equal to the length of the opposite edge augmented by $`2`$. As an edge and its opposite vertex are dual through the morphism $``$ (see Proposition 5.3) and its analog $`^{}`$ attached to the cone $`\sigma ^{}`$, the triangles appearing in the zigzag diagram are a convenient graphical representation of the duality explained in section 5.1. $``$ When $`l_{}`$ is rational and $`l_+`$ is irrational, we draw a little differently the diagram (see Figure 8). The curves representing $`P(\sigma )`$ and $`P(\sigma ^{})`$ start from points $`V_0`$ and $`V_0^{}`$ of a horizontal line representing the line which contains $`l_{}`$. We represent the integral point $`V_1^{}`$ differently from the points $`V_2^{},V_3^{},\mathrm{}`$, because it may not be a vertex of $`P(\sigma ^{})`$, as explained in Proposition 5.3. The length of $`[V_0^{}V_1^{}]`$ is always $`1`$. The relation between the length of an edge and the weight of the opposite vertex is the same as before, with the exception of the triangle $`V_1^{}V_0V_1`$, where the weight of $`V_1^{}`$ is equal to $`l_𝐙[V_0V_1]+1`$. $``$ When both $`l_{}`$ and $`l_+`$ are rational and there is at least one vertex on $`P(\sigma )`$ lying strictly between $`A_{}`$ and $`A_+`$ (that is, $`s1`$), the curves representing $`P(\sigma )`$ and $`P(\sigma ^{})`$ start again from a horizontal line, but now they join in a point $`A_+`$ (see Figure 9). $``$ When both $`l_{}`$ and $`l_+`$ are rational and $`[A_{}A_+]`$ is an edge of $`P(\sigma )`$ (that is, $`s=0`$), the diagram is represented in Figure 10. To summarize, we have the following procedure for constructing and decorating the diagram when $`l_{}`$ is rational: Procedure: Suppose that $`l_{}`$ is rational. Then draw a horizontal line with three marked points $`V_0^{}=A_{}^{},O,V_0=A_{}`$ in this order, $`V_0^{}`$ on the left and $`V_0`$ on the right. Starting from $`V_0^{}`$ and $`V_0`$, draw in the upper half-plane two curves $`P(\sigma ^{})`$, respectively $`P(\sigma )`$, concave towards $`0`$ and coming closer and closer from one another. If $`l_+`$ is rational, join them in a point $`A_+`$. Draw a zigzag line starting from $`V_0`$ and going alternatively from $`P(\sigma )`$ to $`P(\sigma ^{})`$. Denote its successive vertices by $`V_1^{},V_1,V_2^{},\mathrm{}`$ and stop at the point $`V_{s+1}^{}`$. Decorate the edges $`V_0^{}V_1^{}`$ and $`V_{s+1}^{}A_+`$ by $`1`$. The other edges and vertices will be decorated using the initial data (discussed in the sequel), by respecting the following rule: Rule: The weight of a vertex is equal to the length of the opposite edge augmented by the number of its vertices distinct from the points $`A_{},A_{}^{},A_+`$. Initial data: If $`\sigma `$ is of type $`\lambda `$, write the HJ-continued fraction expansion of $`\lambda `$ in the form: (18) $$\lambda =[(2)^{m_1},n_1+3,(2)^{m_2},n_2+3,\mathrm{}]^{}$$ Then decorate the edges of $`P(\sigma )`$ with the numbers $`m_1+1,m_2+1,\mathrm{}`$ and the vertices with the numbers $`n_1+3,n_2+3,\mathrm{}`$. ###### Definition 5.9. We call the previous diagram the zigzag diagram associated to the pair $`(L,\sigma )`$ and to the chosen ordering of the edges of $`\sigma `$, or to the number $`\lambda >1`$, where $`(L,\sigma )`$ is of type $`\lambda `$ with respect to this ordering. We denote it by $`ZZ(\lambda )`$. The zigzag diagrams allow one to visualize the relations between Euclidean and Hirzebruch-Jung continued fractions, proved algebraically in section 2. Indeed, one can read the HJ-continued fraction expansion of $`\lambda >1`$ on the right-hand curved line of $`ZZ(\lambda )`$. By Corollary 5.4, we can read the HJ-continued fraction expansion of $`\frac{\lambda }{\lambda 1}`$ on the left-hand curved line $`P(\sigma )`$ of $`ZZ(\lambda )`$. So, by looking at Figure 9, which can be easily constructed from the initial data by respecting the rule, we get: (19) $$\frac{\lambda }{\lambda 1}=[m_1+2,(2)^{n_1},m_2+3,(2)^{n_2},m_3+3,\mathrm{}]^{}$$ which gives a geometric proof of Proposition 2.7. Now, by Klein’s geometric interpretation of E-continued fractions (see section 3), we see that the E-continued fraction expansion of $`\frac{\lambda }{\lambda 1}`$ can be obtained by writing alternatively the integral lengths of the edges of the polygonal lines $`P(\sigma )`$ and $`P(\sigma ^{})[V_0^{}V_1^{}]`$ (indeed, $`\frac{\lambda }{\lambda 1}`$ is the slope of $`l_+`$ in the base $`(OV_0,OV_1^{})`$): (20) $$\frac{\lambda }{\lambda 1}=[m_1+1,n_1+1,m_2+1,n_2+1,m_3+1,\mathrm{}]^+.$$ This proves geometrically Proposition 2.3. In order to read the E-continued fraction expansion of $`\lambda `$ on the diagram, one has to look at $`ZZ(\lambda )`$ from left to right instead of from right to left and draw a new zigzag line starting from $`V_0^{}`$. The important point here is that one has to discuss according to the alternative $`m_1=0`$ or $`m_1>0`$. In the first case, the zigzag line joins $`V_0^{}`$ to $`V_1`$ and $`V_1`$ to $`V_2^{}`$. In the second case, it joins $`V_0^{}`$ to a new point representing $`A_1`$ and $`A_1`$ to $`V_1^{}`$. Compare this with Lemma 2.5. ###### Example 5.10. Take $`\lambda =\frac{11}{7}`$. After computing $`\lambda =[2,3,2,2]^{}`$, we can construct the associated zigzag diagram $`ZZ(\frac{11}{7})`$. We see that the extreme points $`V_1^{},V_2^{}`$ are vertices of $`P(\sigma ^{})`$. One can read on it the results of the Examples 2.4, 2.6, 2.9. If one had starts instead from $`\lambda =\frac{11}{4}=[3,4]^{}`$, the corresponding diagram would be $`ZZ(\frac{11}{4})`$. In this case the extreme points are not vertices of $`P(\sigma ^{})`$, because their weights are equal to $`2`$. ### 5.3. Relation with the dual cone Denote by $`\stackrel{ˇ}{L}:=\mathrm{Hom}(L,𝐙)`$ the dual lattice of $`L`$. Inside the associated vector space $`\stackrel{ˇ}{L}_𝐑`$ lives the dual cone $`\stackrel{ˇ}{\sigma }`$ of $`\sigma `$, defined by: $$\stackrel{ˇ}{\sigma }:=\{\stackrel{ˇ}{u}\stackrel{ˇ}{L}_𝐑\|\stackrel{ˇ}{u}.u0,u\sigma \}.$$ Let $`\omega `$ be the volume form on $`L_𝐑`$ which verifies $`\omega (u_1,u_2)=1`$ for any basis $`(u_1,u_2)`$ of $`L`$ defining the opposite orientation to $`(l_{},l_+)`$. It is a symplectic form, that is, a non-degenerate alternating bilinear form on $`L_𝐑`$. But we prefer to look at it as a morphism (obtained by making interior products with the elements of $`L`$): $$\omega :L\stackrel{ˇ}{L}.$$ ###### Proposition 5.11. The mapping $`\omega `$ realizes an isomorphism between the pairs $`(L,\sigma ^{})`$ and $`(\stackrel{ˇ}{L},\stackrel{ˇ}{\sigma })`$. Proof: Indeed we have: $$\omega ^1(\stackrel{ˇ}{\sigma })=\{uL\|\omega (u)\stackrel{ˇ}{\sigma }\}=\{uL\|\omega (u,v)0,vL\}=\sigma ^{}.$$ While writing the last equality, we used our convention on the orientation of $`\omega `$. Notice that the dual cone $`\stackrel{ˇ}{\sigma }`$ can be defined without the help of any orientation, in contrast with the morphism $`\omega `$. $`\mathrm{}`$ The previous proposition shows that the construction of the polygonal line $`P(\sigma ^{})`$ explained in Proposition 5.3 describes also the polygonal line $`P(\stackrel{ˇ}{\sigma })`$. This observation is crucial when one wants to use zigzag diagrams for understanding computations with invariants of toric surfaces (see next section). It also helps to understand geometrically the duality between the convex polygons $`K(\sigma )`$ and $`K(\stackrel{ˇ}{\sigma })`$ explained in Gonzalez-Sprinberg and in Oda \[59, pages 27-29\]. As Dimitrios Dais kindly informed us after seeing a version of this paper on ArXiv, a better algebraic understanding of that duality is explained in Dais, Haus & Henk \[14, section 3\]. In particular, modulo Proposition 5.11, the Theorem 3.16 in the previous reference leads easily to an algebraic proof of our Proposition 5.3. (Added in proof) Emmanuel Giroux has informed us that he had realized the existence of a duality between supplementary cones (see \[25, section 1.G\]). ## 6. Relations with toric geometry First we introduce elementary notions of toric geometry (see section 6.1). In section 6.2 we explain how to get combinatorially various invariants of a normal affine toric surface and of the corresponding Hirzebruch-Jung analytic surface singularities. In Section 6.3 we explain how to read the combinatorics of the minimal embedded resolution of a plane monomial curve on an associated zigzag diagram. The basics about resolutions of surface singularities needed in order to understand this section are recalled in section 8.1. ### 6.1. Elementary notions of toric geometry For details about toric geometry, general references are the books of Oda and Fulton , as well as the first survey of it by Kempf, Knudson, Mumford & St. Donat . In the previous section, our fundamental object of study was a pair $`(L,\sigma )`$, where $`L`$ is a lattice of rank 2 and $`\sigma `$ is a strictly convex cone in the 2-dimensional vector space $`L_𝐑`$. Suppose now that the lattice $`L`$ has arbitrary finite rank $`d1`$ and that $`\sigma `$ is a strictly convex rational cone in $`L_𝐑`$. The pair $`(L,\sigma )`$ gives rise canonically to an affine algebraic variety: $$𝒵(L,\sigma ):=\text{Spec}𝐂[\stackrel{ˇ}{\sigma }\stackrel{ˇ}{L}].$$ This means that the algebra of regular functions on $`𝒵(L,\sigma )`$ is generated by the monomials whose exponents are elements of the semigroup $`\stackrel{ˇ}{\sigma }\stackrel{ˇ}{L}`$ of integral points in the dual cone of $`\sigma `$. If $`v\stackrel{ˇ}{\sigma }\stackrel{ˇ}{L}`$, we formally write such a monomial as $`X^v`$. One can show that the variety $`𝒵(L,\sigma )`$ is normal (see the definition at the beginning of section 8.1). The closed points of $`𝒵(L,\sigma )`$ are the morphisms of semigroups $`(\stackrel{ˇ}{\sigma }\stackrel{ˇ}{L},+)(𝐂,)`$. Among them, those whose image is contained in $`𝐂^{}`$ form a $`d`$-dimensional algebraic torus $`𝒯_L=\text{Spec}𝐂[\stackrel{ˇ}{L}]`$, that is, a complex algebraic group isomorphic to $`(𝐂^{})^d`$. The elements of $`L`$ correspond to the 1-parameter subgroups of $`𝒯_L`$, that is, the group morphisms $`(𝐂^{},)(𝒯_L,)`$. The action of $`𝒯_L`$ on itself by multiplication extends canonically to an algebraic action on $`𝒵(L,\sigma )`$, such that $`𝒯_L`$ is the unique open orbit. If $`(\overline{L},\overline{\sigma })`$ is a second pair and $`\varphi :\overline{L}L`$ is a morphism such that $`\varphi (\overline{\sigma })\sigma `$, one gets an associated toric morphism: $$\varphi _{}:𝒵(\overline{L},\overline{\sigma })𝒵(L,\sigma )$$ It is birational if and only if $`\varphi `$ realizes an isomorphism between $`\overline{L}`$ and $`L`$. In this case $`\varphi _{}`$ identifies the tori contained inside $`𝒵(\overline{L},\overline{\sigma })`$ and $`𝒵(L,\sigma )`$. In general: ###### Definition 6.1. Given an algebraic torus $`𝒯`$, a toric variety $`𝒵`$ is an algebraic variety containing $`𝒯`$ as a dense Zariski open set and endowed with an action $`𝒯\times 𝒵𝒵`$ which extends the group multiplication of $`𝒯`$. Oda and Fulton study mainly the normal toric varieties. For an introduction to the study of non-necessarily normal toric varieties, one can consult Sturmfels and González Pérez & Teissier . A normal toric variety can be described combinatorially using fans, that is finite families of rational strictly convex cones, closed under the operations of taking faces or intersections. If $`L`$ is a lattice and $``$ is a fan in $`L_𝐑`$, we denote by $`𝒵(L,)`$ the associated normal toric variety. It is obtained by glueing the various affine toric varieties $`𝒵(L,\sigma )`$ when $`\sigma `$ varies among the cones of the fan $``$. As glueing maps, one uses the toric birational maps $`𝒵(\overline{L},\overline{\sigma })𝒵(L,\sigma )`$ induced by the inclusion morphisms $`(L,\overline{\sigma })(L,\sigma )`$, for each pair $`\overline{\sigma }\sigma `$ of cones of $``$. The variety $`𝒵(L,)`$ is smooth if and only if each cone of the fan $``$ is regular, that is, generated by a subset of a basis of the lattice $`L`$. ### 6.2. Toric surfaces We restrict now to the case of surfaces. Consider a 2-dimensional normal toric surface $`𝒵(L,\sigma )`$, where $`\sigma `$ is a strictly convex cone with non-empty interior. There is a unique 0-dimensional orbit $`O`$, whose maximal ideal is generated by the monomials with exponents in the semigroup $`\stackrel{ˇ}{\sigma }\stackrel{ˇ}{L}O`$. The surface is smooth outside $`O`$, and $`O`$ is a smooth point of it if and only if $`\sigma `$ is a regular cone. Supposing that $`\sigma `$ is not regular, we explain how to describe combinatorially the minimal resolution morphism of $`𝒵(L,\sigma )`$ and the effect of blowing-up the point $`O`$. We also give a formula for the embedding dimension of the germ $`(𝒵(L,\sigma ),O)`$, which is a so-called Hirzebruch-Jung singularity. With the notations of section 4, let us subdivide $`\sigma `$ by drawing the half-lines starting from $`O`$ and passing through the points $`A_k,k\{1,\mathrm{},r\}`$. In this way we decompose $`\sigma `$ in a finite number of regular subcones. They form the minimal regular subdivision of $`\sigma `$, in the sense that any subdivision of $`\sigma `$ by regular cones is necessarily a refinement of the preceding one. The family consisting of the 2-dimensional cones in the subdivision, of their edges and of the origin form a fan $`(\sigma )`$. For each such subcone $`\sigma ^{}`$ of $`\sigma `$, there is a canonical birational morphism $`𝒵(L,\sigma ^{})𝒵(L,\sigma )`$, which realizes an isomorphism of the tori. Using these morphisms, one can glue canonically the tori contained in the surfaces $`𝒵(L,\sigma ^{})`$ when $`\sigma ^{}`$ varies, and obtain a new toric surface $`𝒵(L,(\sigma ))`$, endowed with a morphism: $$𝒵(L,(\sigma ))\stackrel{p_\sigma }{}𝒵(L,\sigma )$$ ###### Proposition 6.2. The morphism $`p_\sigma `$ is the minimal resolution of singularities of the surface $`𝒵(L,\sigma )`$. Moreover, its exceptional locus $`E_\sigma `$ is a normal crossings divisor and the dual graph of $`E_\sigma `$ is topologically a segment. Proof: For details, see . Here we outline only the main steps. The morphism $`p_\sigma `$ is proper, birational and realizes an isomorphism over $`𝒵(L,\sigma )O`$. As $`𝒵(L,(\sigma ))`$ is smooth, $`p_\sigma `$ is a a resolution of singularities of $`𝒵(L,\sigma )`$ (see Definition 8.2). There is a canonical bijection between the irreducible components $`E_k`$ of the exceptional divisor $`E_\sigma =p_\sigma ^1(0)`$ and the half-lines $`[OA_k`$, for $`k\{1,\mathrm{},r\}`$. Moreover, $`E_k`$ is a smooth compact rational curve and (21) $$E_k^2=\alpha _k,k\{1,\mathrm{},r\}$$ where the numbers $`\alpha _k`$ were introduced in relation (10). Using the inequality (11), we deduce that no component of $`E_\sigma `$ is exceptional of the first kind (see the comments which follow Definition 8.2). This implies that $`p_\sigma `$ is the minimal resolution of singularities of $`𝒵(L,\sigma )`$. The proposition is proved. $`\mathrm{}`$ Notice that relation (21) gives an intersection-theoretical interpretation of the weights attached through relation (10) to the integral points situated on $`P(\sigma )`$ which are interior to $`\sigma `$. Conversely (see and ): ###### Proposition 6.3. Suppose that a smooth surface $``$ contains a compact normal crossings divisor $`E`$ whose components are smooth rational curves of self-intersection $`2`$ and whose dual graph is topologically a segment. Denote by $`\alpha _1,\mathrm{},\alpha _r`$ the self-intersection numbers read orderly along the segment. Then $`E`$ can be contracted by a map $`p:(,E)(𝒮,0)`$ to a normal surface $`𝒮`$ and the germ $`(𝒮,0)`$ is analytically isomorphic to a germ of the form $`(𝒵(L,\sigma ),O)`$, where $`\sigma `$ is of type $`\lambda :=[\alpha _1,\mathrm{},\alpha _r]^{}`$. This motivates: ###### Definition 6.4. A normal surface singularity $`(𝒮,0)`$ isomorphic to a germ of the form $`(𝒵(L,\sigma ),O)`$ is called a Hirzebruch-Jung singularity. Hirzebruch-Jung singularities can also be defined as cyclic quotient singularities (see and ). They appear naturally in the so-called Hirzebruch-Jung method of studying an arbitrary surface singularity. Namely, one projects the given singularity by a finite morphism on a smooth surface, then one makes an embedded resolution of the discriminant curve and takes the pull-back of the initial surface by this morphism. In this case, the normalization of the new surface has only Hirzebruch-Jung singularities (see Laufer , Lipman , Brieskorn for details and Popescu-Pampu for a generalization to higher dimensions). The proof of Proposition 6.2 shows that the germs $`(𝒵(L,\sigma ),O)`$ and $`(𝒵(\overline{L},\overline{\sigma }),O)`$ are analytically isomorphic if and only if there exists an isomorphism of the lattices $`L`$ and $`\overline{L}`$ sending $`\sigma `$ onto $`\overline{\sigma }`$. The same is true for strictly convex cones in arbitrary dimensions, as proved by González Pérez & Gonzalez-Sprinberg . Previously we had proved this for simplicial cones in . A Hirzebruch-Jung singularity isomorphic to $`(𝒵(L,\sigma ),O)`$ is said to be of type $`𝒜_{p,q}`$, with $`1q<p`$ and $`\text{gcd}(p,q)=1`$ if (using Definition 5.5) the pair $`(L,\sigma )`$ is of type $`\frac{p}{q}`$ with respect to one of the orderings of the sides of $`\sigma `$. Then, by Proposition 4.3, we have $`\frac{p}{q}=[\alpha _1,\mathrm{},\alpha _r]^{}`$. By Proposition 5.6, one has $`𝒜_{p,q}𝒜_{p^{},q^{}}`$ if and only if $`p=p^{}`$ and $`q^{}\{q,\overline{q}\}`$, where $`q\overline{q}1(modp)`$. The singularities of type $`𝒜_{n+1,n}`$ are also called of type $`𝐀_n`$. They are those for which the polygonal line $`P(\sigma )`$ has only one compact edge, as $`\frac{n+1}{n}=[(2)^n]^{}`$ (a case emphasized in Section 5.2), and also the only Hirzebruch-Jung singularities of embedding dimension $`3`$ (more precisely, they can be defined by the equation $`z^{n+1}=xy`$). Indeed: ###### Proposition 6.5. If $`{\displaystyle \frac{p}{q}}=[\alpha _1,\mathrm{},\alpha _r]^{}=[(2)^{m_1},n_1+3,\mathrm{},n_s+3,(2)^{m_{s+1}}]^{}`$, then: $$embdim(𝒜_{p,q})=3+\underset{i=1}{\overset{r}{}}(\alpha _i2)=\mathrm{\hspace{0.25em}3}+s+\underset{k=1}{\overset{s}{}}n_k.$$ Proof: If $`S`$ is a generating system of the semigroup $`\stackrel{ˇ}{L}\stackrel{ˇ}{\sigma }O`$, then the monomials $`(X^v)_{vS}`$ form a generating system of the Zariski cotangent space $`/^2`$ of the germ at the singular point, where $``$ is the maximal ideal of the local algebra of the singularity $`𝒜_{p,q}`$. By taking a minimal generating system, one gets a basis of this cotangent space. But such a minimal generating system is unique, and consists precisely of the integral points of $`P(\stackrel{ˇ}{\sigma })`$ interior to $`\stackrel{ˇ}{\sigma }`$. By Propositions 5.11 and 2.7, we see that this number is as given in the Proposition. $`\mathrm{}`$ Hirzebruch-Jung singularities are particular cases of rational singularities, introduced by M. Artin , in the 60’s (see also ). In , Tjurina proved that the blow-up of a rational surface singularity is a normal surface which has again only rational singularities (see also the comments of Lê \[49, 4.1\]). As any surface can be desingularized by a sequence of blow-ups of its singular points followed by normalizations (Zariski , see also Cossart and the references therein), this shows that a rational singularity can be desingularized by a sequence of blow-ups of closed points. In particular this is true for a Hirzebruch-Jung singularity. As the operation of blow-up is analytically invariant, we can describe the blow-up of $`O`$ in the model surface $`𝒵(L,\sigma )`$. We use notations introduced at the beginning of the proof of Proposition 5.3. ###### Proposition 6.6. Suppose that the cone $`\sigma `$ is not regular. Subdivide it by drawing the half-lines starting from $`O`$ and passing through the points $`A_1,V_1,`$ $`V_2\mathrm{},V_s,A_r`$. Denote by $`_0(\sigma )`$ the fan obtained in this way. Then the natural toric morphism $`𝒵(L,_0(\sigma ))\stackrel{p_0}{}𝒵(L,\sigma )`$ is the blow-up of $`O`$ in $`𝒵(L,\sigma )`$. Proof: A proof is sketched by Lipman in . Here we give more details. Let $`(𝒮,0)`$ be any germ of normal surface. Consider its minimal resolution $`p_{min}:(_{min},E_{min})(𝒮,0)`$ and its exceptional divisor $`E_{min}=_{k=1}^rE_k`$. The divisors $`Z_{k=1}^r𝐙E_k`$ which satisfy $`ZE_k0,k\{1,\mathrm{},r\}`$ form an additive semigroup with a unique minimal element $`Z_{top}`$, called the fundamental cycle of the singularity. It verifies (22) $$Z_{top}\underset{k=1}{\overset{r}{}}E_k$$ for the componentwise order on the set of cycles with integral coefficients. In the case of a rational singularity, Tjurina showed that the divisors $`E_k`$ which appear in the blow-up of $`0`$ on $`𝒮`$ can be characterized using the fundamental cycle: they are precisely those for which $`Z_{top}E_k<0`$. In our case, where $`(𝒮,0)=(𝒵(L,\sigma ),O)`$, Proposition 6.2 shows that $`p_{min}=p_\sigma `$. Using the relations (21) and (22), we see that $`Z_{top}=_{k=1}^rE_k`$. Again using relation (21), we get: $$Z_{top}E_k<0\text{either }k\{1,r\}\text{ or }\alpha _k3.$$ This shows that the components of $`E_\sigma `$ which appear when one blows-up the origin, are precisely those which correspond to the half-lines $`[OA_1,[OV_1,`$ $`[OV_2,\mathrm{},[OV_s,`$ $`[OA_r`$. But the surface obtained by blowing-up the origin is again normal, by Tjurina’s theorem, which shows that it coincides with $`𝒵(L,\sigma )`$. $`\mathrm{}`$ One sees that after the first blow-up, the new surface has only singularities of type $`𝐀_n`$, where $`n`$ varies in a finite set of positive numbers. The singular points are contained in the set of $`0`$-dimensional orbits of the toric surface $`𝒵(L,_0(\sigma ))`$, which in turn correspond bijectively to the 2-dimensional cones of the fan $`_0(\sigma )`$. The germs of the surface at those points are Hirzebruch-Jung singularities of types $`𝐀_{n_0},\mathrm{},𝐀_{n_s}`$, where $`n_0=l_𝐙[A_1V_1],n_1=l_𝐙[V_1V_2],\mathrm{},n_s=l_𝐙[V_sA_r]`$. We have spoken until now of algebraic aspects of Hirzebruch-Jung singularities. We discuss their topology in section 8.3. ### 6.3. Monomial plane curves Suppose that $`(𝒮,0)`$ is a germ of smooth surface and that $`(𝒞,0)(𝒮,0)`$ is a germ of reduced curve. A proper birational morphism $`p:𝒮`$ is called an embedded resolution of the germ $`(𝒞,0)`$ if $``$ is smooth, $`p`$ is an isomorphism above $`𝒮0`$ and the total transform $`p^1(𝒞)`$ of $`𝒞`$ is a divisor with normal crossings on $``$ in a neighborhood of the exceptional divisor $`E:=p^1(0)`$. The difference $`p^1(𝒞)p^1(0)`$ is called the strict transform of $`𝒞`$ by the morphism $`p`$. It is known since the XIX-th century that any germ of plane curve can be resolved in an embedded way by a sequence of blow-ups of points (see Enriques & Chisini , Laufer , Brieskorn & Knörrer ). The combinatorics of the exceptional divisor of the resolution can be determined starting from the Newton-Puiseux exponents of the irreducible components of the curve and from their intersection numbers using E-continued fraction expansions. We explain here how to read the sequence of self-intersection numbers of the components of the exceptional divisor of the minimal embedded resolution of a monomial plane curve by using a zigzag diagram, instead of just doing blindly computations with continued fractions. If $`p,q𝐍^{},\mathrm{\hspace{0.25em}1}q<p`$ and $`gcd(p,q)=1`$, consider the plane curve $`C_{p/q}`$ defined by the equation: (23) $$x^py^q=0$$ It can be parametrized by: (24) $$\{\begin{array}{c}x=t^q\hfill \\ y=t^p\hfill \end{array}$$ As $`p`$ and $`q`$ are relatively prime, one sees that (24) describes the normalization morphism for $`C_{p/q}`$ (see its definition at the beginning of section 8.1). As $`t^p`$ and $`t^q`$ are monomials, one says that $`C_{p/q}`$ is a monomial curve. There is a natural generalization to higher dimensions (see Teissier ). If one identifies the plane $`𝐂^2`$ of coordinates $`(x,y)`$ with the toric surface $`𝒵(L_0,\sigma _0)`$, where $`L_0=𝐙^2`$ and $`\sigma _0`$ is the first quadrant, then it is easy to see (look at equation (24)) that $`C_{p/q}`$ is the closure in $`𝐂^2`$ of the image of the 1-parameter subgroup of the complex torus $`𝒯_{L_0}=(𝐂^{})^2`$ corresponding to the point $`(q,p)`$. Consider again the notations introduced before Lemma 3.1. Let $`l_{}:=[O(1,0)`$ and $`l_+:=[O(q,p)`$ be the edges of the cone $`\sigma _x(\frac{p}{q})`$. We leave to the reader the proof of the following lemma, which is very similar to the proof of Lemma 3.1. Recall that the type of a cone was introduced in Definition 5.5. ###### Lemma 6.7. With respect to the chosen ordering of its edges, the cone $`\sigma _x(\frac{p}{q})`$ is of type $`\frac{p}{pq}`$. Moreover, with the notations of section 5, $`A_1=(1,1)`$, $`A_1^{}=(0,1)`$ and $`A_+=(q,p)`$. Even if the proof is very easy, it is important to be conscious of this result, as it allows to apply the study done in section 5 to our context. Given the pair $`(p,q)`$, we want to describe the process of embedded resolution of the curve $`C_{p/q}`$ by blow-ups, as well as the final exceptional divisor, the self-intersections of its components and their orders of appearance during the process. ###### Lemma 6.8. The blow-up $`\pi _0:_0𝐂^2`$ of $`0`$ in $`𝐂^2`$ is a toric morphism corresponding to the subdivision of $`\sigma _0`$ obtained by joining $`O`$ to $`A_1=(1,1)`$. The strict transform of $`C_{p/q}`$ passes through the $`0`$-dimensional orbit of $`_0`$ associated to the cone $`𝐑_+OA_1+𝐑_+OA_1^{}`$. Proof: With the notations of Section 3, we consider the fan $`_0`$ subdividing $`\sigma _0`$ which consists of the cones $`\sigma _x(1),\sigma _y(1)`$, their edges and the origin. Let $`\pi __0:𝒵(L,_0)𝒵(L,\sigma _0)`$ be the associated toric morphism. It is obtained by gluing the maps $`\pi _x:𝒵(L,\sigma _x(1))𝒵(L,\sigma _0)`$ and $`\pi _y:𝒵(L,\sigma _y(1))𝒵(L,\sigma _0)`$ over $`(𝐂^{})^2`$. With respect to the coordinates given by the monomials associated to the primitive vectors of $`L`$ situated on the edges of the cones $`\sigma _0,\sigma _x(1),\sigma _y(1)`$, the maps $`\pi _x`$ and $`\pi _y`$ are respectively described by: $$\{\begin{array}{c}x=x_1y_1\hfill \\ y=y_1\hfill \end{array}\mathrm{and}\{\begin{array}{c}x=x_2\hfill \\ y=x_2y_2\hfill \end{array}$$ One recognizes the blow-up of $`0`$ in $`𝐂^2`$. Now, in order to compute the strict transform of $`C_{p/q}`$, one has to make the previous changes of variables in equation (19). The lemma follows immediately. $`\mathrm{}`$ Starting from Lemma 6.7 and using the previous lemma as an induction step, we get: ###### Proposition 6.9. The following procedure constructs the dual graph of the total transform of $`C_{p/q}`$ by the minimal embedded resolution morphism, starting from the zigzag diagram $`ZZ(\frac{p}{pq})`$: $``$ On each edge of integral length $`l1`$, add $`(l1)`$ vertices of weight $`2`$. Then erase the weights of the edges (that is, their integral length). $``$ Attach the weight $`1`$ to the vertex $`A_+`$. Then change the signs of all the weights of the vertices. $``$ Label the vertices by the symbols $`E_1,E_2,E_3,\mathrm{}`$ starting from $`A_1`$ on $`P(\sigma )`$ till arriving at $`V_1`$, continuing from the first vertex which follows $`V_1^{}`$ on $`P(\sigma ^{})`$ till arriving at $`V_2^{}`$, coming then back to $`P(\sigma )`$ at the first vertex which follows $`V_1`$ and so on, till labelling the vertex $`A_+`$. $``$ Erase the horizontal line, the zigzag line and the curved segment between $`V_0^{}`$ and the first vertex which follows $`V_1^{}`$. $``$ Add an arrow to the vertex $`A_+`$ and keep only the weights of the vertices and their labels $`E_n`$. The arrowhead vertex represents the strict transform of the curve $`C_{p/q}`$ and the indices of the components $`E_i`$ correspond to the orders of appearance during the process of blow-ups. It is essential to remark that in the previous construction one starts from $`ZZ(\frac{p}{pq})`$ and not from $`ZZ(\frac{p}{q})`$ (look again at Lemma 6.7). ###### Example 6.10. Consider the curve $`x^{11}y^4=0`$. Then $`\lambda =\frac{11}{114}=\frac{11}{7}`$. Its zigzag diagram $`ZZ(\frac{11}{7})`$ was constructed in Example 5.10. So, the dual graph of the total transform of $`C_{11/4}`$ by the minimal embedded resolution morphism has 6 vertices, of easy computable weights (see Figure 13). Proposition 6.9 endows us with an easy way of remembering the following classical description of the minimal embedded resolution of a monomial plane curve (see Jurkiewicz , who attributes it to Hirzebruch; Spivakovsky extends it to the case of monomial-type valuations on function-fields of surfaces): ###### Proposition 6.11. If $`{\displaystyle \frac{p}{q}}=[m_1+1,n_1+1,m_2+1,\mathrm{},n_s+1,m_{s+1}+1]^+`$, then the dual graph of the total transform of the monomial curve $`C_{p/q}`$ is the one which appears in Figure 14. Proof: Combine formulae (20) and (18) with Figure 9 and Proposition 6.9. $`\mathrm{}`$ In Figure 14 we have indicated only the orders of appearance of the components of the exceptional divisor corresponding to the extremities of the graph. We leave as an exercise for the reader to complete the diagram with the sequence $`(E_k)_{k1}`$. Notice that in the E-continued fraction expansion of $`\frac{p}{q}`$ used in the previous proposition, there is the possibility that $`m_{s+1}=0`$. In this case, the canonical expansion is obtained using relation (7). But in order to express in a unified form the result of the application of the algorithm, it was important for us to use an expansion of $`\frac{p}{q}`$ with an odd number of partial quotients (which is always possible, precisely according to formula (7)). One can use the combinatorics of the embedded resolution of monomial plane curves as building blocks for the description of the combinatorics of the resolution of any germ of plane curve. A detailed description of the passage between the Eggers tree, which encodes the Newton-Puiseux exponents of the components of the curve, and the dual graph of the total transform of the curve by its embedded resolution morphism can be found in García Barroso (see also Brieskorn & Knörrer \[9, section 8.4\] and Wall ). A topological interpretation of the trees appearing in these two encodings was given in Popescu-Pampu \[61, chapter 4\]. In higher dimensions, González Pérez used toric geometry in order to describe embedded resolutions of quasi-ordinary hypersurface singularities. Again, the building blocks are monomial varieties. A prototype for his study is the method of resolution of an irreducible germ of plane curve by only one toric morphism, developed by Goldin & Teissier . In the classical treatise of Enriques & Chisini , resolutions of curves by blow-ups of points are not studied using combinatorics of divisors, but instead using the infinitely near points through which the strict transforms of the curve pass during the process of blowing ups. Those combinatorics were also encoded in a diagram, called nowadays Enriques diagram (see Casas-Alvero ). Enriques diagrams are very easily constructed using the knowledge of the orders of appearance of the divisors during the process of blowing ups. For this reason, zigzag diagrams combined with Proposition 6.9 give an easy way to draw them for a monomial plane curve. We leave the details to the interested reader. Then one uses this again as building blocks for the analysis of general plane curve singularities (see ). ## 7. Graph structures and plumbing structures on 3-manifolds This section contains preparatory material for the topological study of the 3-manifolds appearing as abstract boundaries of normal surface singularities, done in sections 8 and 9. We recall general facts about Seifert, graph and plumbing structures on 3-manifolds, as well as about JSJ theory. We also define particular classes of plumbing structures on thick tori and solid tori, starting from naturally arising pairs $`(L,\sigma )`$, where $`L`$ is a 2-dimensional lattice and $`\sigma `$ is a rational strictly convex cone in $`L_𝐑`$. Namely, given a pair of essential curves on the boundary of a thick torus $`M`$, their classes generate two lines in the lattice $`L:=H_1(M,𝐙)`$. A choice of orientations of these lines distinguishes one of the four cones in which the lines divide the plane… ### 7.1. Generalities on manifolds and their splittings We denote by $`𝐈`$ the interval $`[0,1]`$, by $`𝐃`$ the closed disc of dimension $`2`$ and by $`\text{S}^n`$ the sphere of dimension $`n`$. An annulus is a surface diffeomorphic to $`I\times 𝐒^1`$. A simple closed curve on a 2-dimensional torus is called essential if it is non-contractible. It is classical that an oriented essential curve on a torus $`T`$ is determined up to isotopy by its image in $`H_1(T,𝐙)`$ (see \[21, section 2.3\]). Moreover, the vectors of $`H_1(T,𝐙)`$ which are homology classes of essential curves are precisely the primitive ones. We say that a manifold is closed if it is compact and without boundary. If $`M`$ is a manifold with boundary, we denote by $`\stackrel{}{M}`$ its interior and by $`M`$ its boundary. If moreover $`M`$ is oriented, we orient $`M`$ in such a way that at a point of $`M`$, an outward pointing tangent vector to $`M`$, followed by a basis of the tangent space to $`M`$, gives a basis of the tangent space to $`M`$ (this is the convention which makes Stokes’ theorem $`_M𝑑\omega =_M\omega `$ true). We say then that $`M`$ is oriented compatibly with $`M`$. If $`M`$ is an oriented manifold, we denote by $`M`$ the same manifold with reversed orientation. If $`M`$ is a closed oriented surface, then $`M`$ is orientation-preserving diffeomorphic to $`M`$. This fact is no longer true in dimension 3, that is why it is important to describe carefully the choice of orientation. In this sense, see Theorem 8.11, as well as Propositions 9.3 and 9.6. We denote by $`\mathrm{Diff}(M)`$ the group of self-diffeomorphisms of $`M`$, by $`\mathrm{Diff}^{}(M)`$ the subgroup of self-diffeomorphisms which are isotopic to the identity and by $`\mathrm{Diff}^+(M)`$ the subgroup of diffeomorphisms which preserve the orientation of $`M`$ (when $`M`$ is orientable). ###### Definition 7.1. Let $`M`$ be a 3-manifold with boundary. We say that $`M`$ is a thick torus if it is diffeomorphic to $`𝐒^1\times 𝐒^1\times 𝐈`$. We say that $`M`$ is a solid torus if it is diffeomorphic to $`𝐃\times 𝐒^1`$. We say that $`M`$ is a thick Klein bottle if it is diffeomorphic to a unit tangent circle bundle to the Möbius band. In the definition of a thick Klein bottle $`M`$ we use an arbitrary riemannian metric on a Möbius band. The manifold obtained like this is independent of the choices up to diffeomorphism. Moreover, it is orientable, because any tangent bundle is orientable and the manifold we define appears as the boundary of a unit tangent disc bundle. The preimage of a central circle of the Möbius band by the fibration map is a Klein bottle, and the manifold $`M`$ appears then as a tubular neighborhood of it, which explains the name. For details, see \[81, section 3\] and \[21, section 10.11\]. On the boundary of a solid torus $`M`$ there exists an essential curve which is contractible in $`M`$. Such a curve, which is unique up to isotopy (see ), is called a meridian of $`M`$. A 3-manifold $`M`$ is called irreducible if any embedded sphere bounds a ball. A surface embedded in $`M`$ is called incompressible if its $`\pi _1`$ injects in $`\pi _1(M)`$. Two tori embedded in $`M`$ are called parallel if they are disjoint and they cobound a thick torus embedded in $`M`$. The manifold $`M`$ is called atoroidal if any embedded incompressible torus is parallel to a component of $`M`$. ###### Definition 7.2. Let $`M`$ be an orientable manifold and $`S`$ be an orientable closed (not necessarily connected) hypersurface of $`M`$. A manifold with boundary $`M_S`$ endowed with a map $`M_S\stackrel{r_{M,S}}{}M`$ is called a splitting of $`\mathrm{M}`$ along $`\mathrm{S}`$ if: $``$ $`r_{M,S}`$ is a local embedding; $``$ $`M_S=(r_{M,S})^1(S)`$ and the restriction $`r_{M,S}\|_{M_S}`$ is a trivial double covering of $`S`$; $``$ the restriction $`(r_{M,S})\|_{\underset{S}{\overset{}{M}}}:\underset{S}{\overset{}{M}}MS`$ is a diffeomorphism. If this is the case, the map $`r_{M,S}`$ is called the reconstruction map associated to the splitting. We say that $`S`$ splits $`M`$ into $`M_S`$ and that the connected components of $`M_S`$ are the pieces of the splitting. If $`N`$ is a piece of $`M_S`$ and $`PM`$ is a set, we say that $`P`$ contains $`N`$ if $`r_{M,S}(N)P`$. It can be shown easily that splittings of $`M`$ along $`S`$ exist and are unique up to unique isomorphism. The idea is very intuitive, one simply thinks at $`M`$ being split open along each connected component of $`S`$. A way to realize this is to take the complement of an open tubular neighborhood of $`S`$ in $`M`$ and to deform the inclusion mapping in an arbitrarily small neighborhood of the boundary in order to push it towards $`S`$ (see Waldhausen and Jaco ). If $`\varphi \mathrm{Diff}^+(M)`$, one can also canonically split $`\varphi `$ and get a diffeomorphism $`\varphi _S`$ of manifolds with boundary (we leave the axiomatic definition of $`\varphi _S`$ to the reader): $$\varphi _S:M_SM_{\varphi (S)}$$ Among closed 3-manifolds, two particular classes will be especially important for us, the lens spaces and the torus fibrations. The reason why we treat them simultaneously will appear clearly in section 8.3. ###### Definition 7.3. Let $`M`$ be an orientable 3-manifold. We say that $`M`$ is a lens space if it contains an embedded torus $`T`$ such that $`M_T`$ is the disjoint union of two solid tori whose meridians have non-isotopic images on $`T`$. We say that $`M`$ is a torus fibration if it contains an embedded torus $`T`$ such that $`M_T`$ is a thick torus. Lens spaces can also be defined as quotients of $`𝐒^3`$ by linear free cyclic actions or - and this explains the name - as manifolds obtained by gluing in a special way the faces of a lens-shaped polyhedron (see or \[21, section 4.3\]). We impose the condition on the meridians in order to avoid the manifold $`𝐒^1\times 𝐒^2`$, which can also be split into two solid tori, but whose universal cover is not the 3-dimensional sphere, a difference which makes it to be excluded from the set of lens spaces by most authors. There exists a classical encoding of oriented lens spaces by positive integers. We recall it at the end of section 9.1 (see Proposition 9.4). If $`M`$ is a torus fibration and $`TM`$ splits it into a thick torus, then a trivial foliation of $`M_T`$ by tori parallel to the boundary components is projected by $`r_{M,T}`$ onto a foliation by pairwise parallel tori. The space of leaves is topologically a circle and the projection $`\pi :M𝐒^1`$ is a locally trivial fibre bundle whose fibres are tori, which explains the name. ###### Definition 7.4. Let $`\pi :M𝐒^1`$ be a locally trivial fibre bundle whose fibres are tori. Fix a fibre of $`\pi `$ (for example the initial torus $`T`$) and also an orientation of the base space $`𝐒^1`$. The algebraic monodromy operator $`m`$ is by definition the first return map of the natural parallel transport on the first homology fibration over $`𝐒^1`$, when one travels in the positive direction. The map $`m`$ is a well-defined linear automorphism $`mSL(H_1(T,𝐙))`$, once an orientation of $`𝐒^1`$ was chosen. Its conjugacy class in $`SL(2,𝐙)`$ is independent of the choice of the fibre. If one changes the orientation of $`𝐒^1`$, then $`m`$ is replaced by $`m^1`$. This shows that the trace of $`m`$ is independent of the choice of $`T`$ and of the orientation of $`𝐒^1`$. Remark that no choice of orientation of $`M`$ is needed in order to define it. For more information about torus fibrations, see Neumann and Hatcher . We come back to them in Section 9.2, with special emphasis on subtleties related to their orientations. ### 7.2. Seifert structures Seifert manifolds are special 3-manifolds whose study can be reduced in some way to the study of lower-dimensional spaces. ###### Definition 7.5. A Seifert structure on a 3-manifold $`M`$ is a foliation by circles such that any leaf has a compact orientable saturated neighborhood. A leaf with trivial holonomy is called a regular fibre. A leaf which is not regular is called an exceptional fibre. The space of leaves is called the base of the Seifert structure. We say that a Seifert structure is orientable if there is a continuous orientation of all the leaves of the foliation. If such an orientation is fixed, one says that the Seifert structure is oriented. If there exists a Seifert structure on $`M`$, we say that $`M`$ is a Seifert manifold. The condition on the leaves to have compact saturated neighborhoods is superfluous if the ambient manifold $`M`$ is compact, it is enough then to ask that any leaf be orientation-preserving, as was shown by Epstein . This is no longer true on non-compact manifolds, as was shown by Vogt . The initial definition of Seifert was slightly different: a) He did not speak of “foliation”, but of “fibration”. b) He gave models for the possible neighborhoods of the leaves. In what concerns point a), Seifert’s definition is one of the historical sources of the concept of fibration and fibre bundle. For him a fibration is a decomposition of a manifold into “fibres”; only in a second phase can one try to construct the associated “orbit space”, or the “base” with our vocabulary. This shows that his definition is closer to the present notion of foliation; in fact his “fibration” is a foliation, but this can be seen only by using the required condition on model neighborhoods. We prefer to speak about “Seifert structure” and not “Seifert fibration” precisely because what is important to us is to see the structure as living inside the manifold, which makes possible to speak about isotopies. For details about the historical development of different notions of fibrations, see Zisman . In what concerns point b), the possible orientable saturated neighborhoods of foliations by circles coincide up to a leaf-preserving diffeomorphism with Seifert’s model neighborhoods. If one drops the orientability condition, appears a new model which was not considered by Seifert, but which is very useful in the classification of non-orientable 3-manifolds (see Scott , Bonahon ). Some general references about Seifert manifolds are Orlik , Neumann & Raymond (where the base was defined as an orbifold), Scott , Fomenko & Matveev and Bonahon . In the sequel, we are interested in Seifert structures only up to isotopy. ###### Definition 7.6. Two Seifert structures $`_1`$ and $`_2`$ on $`M`$ are called isotopic if there exists $`\varphi \mathrm{Diff}^{}(M)`$ such that $`\varphi (_1)=_2`$. The following proposition is proved in Jaco and Fomenko & Matveev . ###### Proposition 7.7. The only orientable compact connected 3-manifolds with non-empty boundary which admit more than one Seifert structure up to isotopy are the thick torus, the solid torus and the thick Klein bottle. a) If $`M`$ is a thick torus, any essential curve on one of its boundary components is the fibre of a Seifert structure on $`M`$, unique up to isotopy, and devoid of exceptional fibres. Moreover, $`M`$ appears like this as the total space of a trivial circle bundle over an annulus. b) If $`M`$ is a solid torus and $`\gamma `$ is a meridian of it, an essential curve $`c`$ on its boundary is a fibre of a Seifert structure on $`M`$ if and only if their homological intersection number $`[c][\gamma ]`$ (once they are arbitrarily oriented) is non-zero. In this case, the associated structure is unique up to isotopy and has at most one exceptional fibre. All fibres are regular if and only if $`[c][\gamma ]=\pm 1`$. In this last case, $`M`$ appears as the total space of a trivial circle bundle over a disc. c) If $`M`$ is a thick Klein bottle, it admits up to isotopy two Seifert structures. One of them is devoid of exceptional fibres and its space of orbits is a Möbius band. The other one has two exceptional fibres with holonomy of order $`2`$ and its space of orbits is topologically a disc. The closed orientable 3-manifolds which admit more than one Seifert structure up to isotopy are also classified (see Bonahon and the references therein). In this paper we need only the following less general result, which can be deduced by combining with (see Definition 8.1): ###### Proposition 7.8. The only 3-manifolds which are diffeomorphic to abstract boundaries of normal surface singularities and which admit non-isotopic Seifert structures are the lens spaces. ### 7.3. Graph structures and JSJ decomposition theory If one glues various Seifert manifolds along components of their boundaries, one obtains so-called graph-manifolds: ###### Definition 7.9. A graph structure on a 3-manifold $`M`$ is a pair $`(𝒯,)`$, where $`𝒯`$ is an embedded surface in $`M`$ whose connected components are tori and where $``$ is a Seifert structure on $`M_𝒯`$ (see Definition 7.2). We say that a graph structure is orientable if $``$ is an orientable Seifert structure on $`M_𝒯`$. If there exists a graph structure on $`M`$, we say that $`M`$ is a graph manifold. Notice that no particular graph structure is specified when one speaks about a graph manifold. One only supposes that there exists one. In the sequel we are interested in graph structures on a given manifold only up to isotopy: ###### Definition 7.10. Two graph structures $`(𝒯_1,_1)`$, $`(𝒯_2,_2)`$ on $`M`$ are called isotopic if there exists $`\varphi \mathrm{Diff}^{}(M)`$ such that $`\varphi (𝒯_1)=𝒯_2`$ and $`\varphi _{𝒯_1}(_1)`$ is isotopic to $`_2`$. Graph manifolds were introduced by Waldhausen , generalizing von Randow’s tree manifolds (see their definition in the next paragraph) studied in . Following Mumford who proved Poincaré conjecture for the abstract boundaries of normal surface singularities (see Definition 8.1), von Randow proved it for tree manifolds; his proof contained a gap which was later filled by Scharf . Waldhausen’s definition was different from Definition 7.9. On one side he did not allow exceptional fibres in the Seifert structure on $`M_𝒯`$ and on another side he did not fix (up to isotopy) a precise fibration by circles, but only supposed that such a fibration existed. He represented a graph structure by a finite graph with decorated vertices and edges (corresponding respectively to the pieces of $`M_𝒯`$ and to the components of $`𝒯`$), which explains the name. Tree manifolds are then the graph manifolds which admit a graph structure $`(𝒯,)`$ such that the corresponding graph is a tree and the base of the Seifert structure on $``$ has genus $`0`$. With our definition, graph structures can also be encoded by graphs. One has only to add more decorations to the vertices, in order to keep in memory the exceptional fibres of the corresponding Seifert fibred pieces. With his definition, Waldhausen solved the homeomorphism problem for graph-manifolds, by giving normal forms for the graph structures on a given manifold and by showing that with exceptions in a finite explicit list, any irreducible graph-manifold has a graph-structure in normal form which is unique up to isotopy. Later, Jaco & Shalen and Johannson showed that there remains no exception in the classification up to isotopy if one modifies the notion of graph-structure by allowing exceptional fibres, that is, when one works with Definition 7.9. More generally, they proved: ###### Theorem 7.11. Let $`M`$ be a compact, connected, orientable and irreducible $`3`$-manifold (with possible non-empty boundary). Then $`M`$ contains an embedded surface $`𝒯`$ whose connected components are incompressible tori and such that any piece of $`M_𝒯`$ is either a Seifert manifold or is atoroidal. Moreover, if $`𝒯`$ is minimal for the inclusion among surfaces with this property, then it is well-defined up to isotopy. We say that a minimal family $`𝒯`$ as in the previous theorem is a JSJ family of tori in $`M`$. A variant of the previous theorem considers also embedded annuli. These various theorems of canonical decomposition are called nowadays Jaco-Shalen-Johannson (JSJ) decomposition theory, and were the starting point of Thurston’s geometrization program, as well as of the theory of JSJ decompositions for groups. For details about JSJ decompositions, in addition to the previously quoted books one can consult Jaco , Neumann & Swarup , Hatcher and Bonahon . In and , we showed that also knot theory inside an irreducible 3-manifold reflects the ambient JSJ decomposition. We define now a notion of minimality for graph structures on a given manifold. ###### Definition 7.12. Suppose that $`(𝒯,)`$ is a graph structure on $`M`$. We say that it is minimal if the following conditions are verified: $``$ No piece of $`M_𝒯`$ is a thick torus or a solid torus. $``$ One cannot find a Seifert structure $`^{}`$ on $`M_𝒯`$ such that the images of its leaves by the reconstruction mapping $`r_{M,𝒯}`$ coincide on a component of $`𝒯`$. As a corollary of Theorem 7.11, if $`(𝒯,)`$ is a minimal graph structure on $`M`$, then $`𝒯`$ is the minimal JSJ system of tori in $`M`$. But one can prove more: ###### Theorem 7.13. Each closed orientable irreducible graph manifold which is not a torus fibration with $`\|trm\|3`$ admits a minimal graph structure. Moreover, the family $`𝒯`$ of tori associated to a minimal graph structure coincides with the JSJ family of tori. In particular, it is unique up to an isotopy. Suppose that $`(𝒯,)`$ is a given graph structure without thick tori and solid tori among its pieces. In view of Proposition 7.7, its only pieces which can have non-isotopic Seifert structures are the thick Klein bottles. This shows that, in order to check whether $`(𝒯,)`$ is minimal or not, one has only to consider the possible choices of Seifert structures on them up to isotopy (that is $`2^n`$ possibilities, where $`n`$ is the number of such pieces). Suppose that $`M`$ is a graph manifold which is neither a torus fibration with $`\|trm\|3`$, nor a Seifert manifold which admits non-isotopic Seifert structures. Then, if $`𝒯`$ is a family of tori associated to a minimal graph structure, there is a unique Seifert structure on $`M_𝒯`$ up to isotopy, such that each piece which is a thick Klein bottle has an orientable base. ###### Definition 7.14. Suppose that $`M`$ is an orientable graph manifold which is neither a torus fibration with $`\|trm\|3`$ nor a Seifert manifold which admits non-isotopic Seifert structures. We say that a minimal graph structure is the canonical graph structure on $`M`$ if each piece which is a thick Klein bottle has an orientable base. ### 7.4. Plumbing structures Plumbing structures are special types of graph structures: ###### Definition 7.15. A plumbing structure on a 3-manifold $`M`$ is a graph structure without exceptional fibres $`(𝒯,)`$ on $`M`$, such that for any component $`T`$ of $`𝒯`$, the homological intersection number on $`T`$ of two fibres of $``$ coming from opposite sides is equal to $`\pm 1`$. Plumbing structures are the ancestors of graph structures. They were introduced by Mumford in the study of singularities of complex analytic surfaces (see Hirzebruch , Hirzebruch, Neumann & Koh , as well as our explanations in section 8.2). In fact Mumford does not speak about “plumbing structure”. Instead, he describes a way to construct the abstract boundary of a normal surface singularity (see Definition 8.1) by gluing total spaces of circle-bundles over real surfaces using “plumbing fixtures”. Later on, “plumbing” was more used as a verb than as a noun. That is, one concentrated more on the operations needed to construct a new object from elementary pieces, than on the structure obtained on the manifold resulting from the construction. The fact that we are interested precisely in this structure up to isotopy and not on the graph which encodes it, is a difference with Neumann for example. In , Neumann describes an algorithm for deciding if two manifolds obtained by plumbing are diffeomorphic. He uses as an essential ingredient Waldhausen’s classification theorem of graph manifolds (according to the definition which does not allow exceptional fibres, see the comments made in section 7.3). In fact, by using the uniqueness up to isotopy of the JSJ-tori, we can deduce the uniqueness up to isotopy for special plumbing structures on singularity boundaries. This is the subject of section 9. Even if Definition 7.15 seems to suggest the opposite, the class of graph manifolds is the same as the class of manifolds which admit a plumbing structure. A way to see this is to use the construction of plumbing structures on thick tori and solid tori described in section 7.5. For a detailed comparison of graph structures and plumbing structures, as well as for a study of the elementary operations on them, one can consult Popescu-Pampu \[61, chapter 4\]. ### 7.5. Hirzebruch-Jung plumbing structures on thick tori and solid tori In this section we define special classes of plumbing structures on thick tori and solid tori, which will be used in section 9. The starting point is in both cases a pair $`(L,\sigma )`$ of a 2-dimensional lattice and a rational strictly convex cone $`\sigma L_𝐑`$, naturally attached to essential curves on the boundary of the 3-manifold. $``$ Suppose first that $`M`$ is an oriented thick torus. On each component of its boundary, we consider an essential curve. Denote by $`\gamma ,\delta `$ these curves. We suppose that their homology classes (once they are arbitrarily oriented) in $`H_1(M,𝐑)𝐑^2`$ are non-proportional. So, we are in presence of a 2-dimensional lattice $`L=H_1(M,𝐙)`$ and of two distinct rational lines in it, generated by the homology classes $`[\gamma ],[\delta ]`$. Orient $`M`$ compatibly with $`M`$. Then order in an arbitrary way the components of $`M`$: call the first one $`T_{}`$ and the second one $`T_+`$. Denote by $`\gamma _{}`$ the simple closed curve drawn on $`T_{}`$ and by $`\gamma _+`$ the one drawn on $`T_+`$. Then orient $`\gamma _{}`$ and $`\gamma _+`$. By hypothesis, their homology classes $`[\gamma _{}],[\gamma _+]`$ are non-proportional primitive vectors in the 2-dimensional lattice $`L=H_1(M,𝐙)`$. This shows that $`([\gamma _{}],[\gamma _+])`$ is a basis of $`L_𝐑=H_1(M,𝐑)`$ which induces an orientation of this vector space. As $`T_+`$ is a deformation retract of $`M`$, one has canonically $`H_1(T_+,𝐙)=L`$, and so the ordered pair $`(\gamma _{},\gamma _+)`$ induces an orientation of $`T_+`$. ###### Definition 7.16. We say that $`\gamma _{}`$ and $`\gamma _+`$ are oriented compatibly with the orientation of $`M`$ if, when taken in the order $`(\gamma _{},\gamma _+)`$, they induce on $`T_+`$ an orientation which coincides with its orientation as a component of $`M`$. Of course, a priori there is no reason for choosing this notion of compatibility rather than the opposite one. Our choice was done in order to get a more pleasant formulation for Lemma 8.5. Let $`\sigma `$ be the cone generated by $`[\gamma _{}]`$ and $`[\gamma _+]`$ in $`L_𝐑`$. As these homology classes were supposed non-proportional, the cone $`\sigma `$ is strictly convex and has non-empty interior. Denote by $`l_\pm `$ the edge of $`\sigma `$ which contains the integral point $`[\gamma _\pm ]`$. Then, with the notations of section 4, $`A_\pm =[\gamma _\pm ]`$. Indeed, as $`\gamma _\pm `$ is an essential curve of $`T_\pm `$, its homology class is a primitive vector of $`L`$. Let $`(A_n)_{0nr+1}`$ be the integral points on the compact edges of $`P(\sigma )`$, defined in section 4. So, $`OA_0=[\gamma _{}]`$ and $`OA_{r+1}=[\gamma _+]`$. Let $`(T_n)_{0nr+2}`$ be a sequence of pairwise parallel tori in $`M`$, such that $`T_0=T_{}`$ and $`T_{r+2}=T_+`$. Moreover, we number them in the order in which they appear between $`T_{}`$ and $`T_+`$. Denote $`𝒯:=_{n=1}^{r+1}T_n`$. If $`M_n`$ denotes the piece of $`M_𝒯`$ whose boundary components are $`T_n`$ and $`T_{n+1}`$, where $`n\{0,\mathrm{},r+1\}`$, we consider on it a Seifert structure such that the homology class of its fibres in $`L`$ is $`OA_n`$. We get like this a plumbing structure on $`M`$, well-defined up to isotopy, and depending only on the triple $`(M,\gamma _{},\gamma _+)`$. We see that the simultaneous change of the orientations of $`\gamma _{}`$ and $`\gamma _+`$ or the change of their ordering (in order to respect the compatibility condition of Definition 7.16) leads to the same (unoriented) plumbing structure. ###### Definition 7.17. We say that the previous unoriented plumbing structure on the oriented thick torus $`M`$ is the Hirzebruch-Jung plumbing structure associated to $`(\gamma ,\delta )`$ and we denote it by $`𝒫(M,\gamma ,\delta )`$. $``$ Suppose now that $`M`$ is an oriented solid torus. We consider an essential curve $`\gamma `$ on $`M`$ which is not a meridian. Take a torus $`T`$ embedded in $`\stackrel{}{M}`$ and parallel to $`M`$. Denote by $`N`$ the thick torus contained between $`M`$ and $`T`$. Put $`T_{}=M,T_+=T,\gamma _{}=\gamma `$ and let $`\gamma _+`$ be an essential curve on $`T_+`$ which is a meridian of the solid torus $`M\stackrel{}{N}`$ (see Figure 16). Consider the Hirzebruch-Jung plumbing structure $`𝒫(N,\gamma _{},\gamma _+)`$. With the notations of the construction done for thick tori, denote $`𝒯(M,\gamma ):=_{n=1}^rT_n`$. Then the pieces of $`M_{𝒯(M,\gamma )}`$ are the thick tori $`M_0,M_1,\mathrm{},M_{r1}`$ and a solid torus which is the “union” of $`M_r,M_{r+1}`$ and $`M\stackrel{}{N}`$. On $`M_0,\mathrm{},M_{r1}`$ we keep the Seifert structure of $`𝒫(N,\gamma _{},\gamma _+)`$. On the solid torus we extend the Seifert structure of $`M_r`$. By Proposition 7.7 b), we see that this Seifert structure has no exceptional fibres. This shows that we have constructed a plumbing structure on $`M`$. It is obviously well-defined up to isotopy, once the isotopy class of $`\gamma `$ is fixed. ###### Definition 7.18. We say that the previous unoriented plumbing structure on the oriented solid torus $`M`$ is the Hirzebruch-Jung plumbing structure associated to $`\gamma `$ and we denote it by $`𝒫(M,\gamma )`$. ## 8. Generalities on the topology of surface singularities In this section we look at the boundaries $`M(𝒮)`$ of normal surface singularities $`(𝒮,0)`$. We explain how to associate to any normal crossings resolution $`p`$ of $`(𝒮,0)`$ a plumbing structure on $`M(𝒮)`$. Then we explain how to pass from the plumbing structure associated to the minimal normal crossings resolution of $`(𝒮,0)`$ to the canonical graph structure on $`M(𝒮)`$ (see Definition 7.14). We recommend the survey articles of Némethi and Wall for an introduction to the classification of normal surface singularities. ### 8.1. Resolutions of normal surface singularities and their dual graphs First we recall basic facts about normal analytic spaces. Let $`𝒱`$ be a reduced analytic space. It is called normal if for any point $`P𝒱`$, the germ $`(𝒱,P)`$ is irreducible and its local algebra is integrally closed in its field of fractions. If $`𝒱`$ is not normal, then there exists a finite map $`\nu :\stackrel{~}{𝒱}𝒱`$ which is an isomorphism over a dense open set of $`𝒱`$ and such that $`\stackrel{~}{𝒱}`$ is normal. Such a map, which is unique up to unique isomorphism, is called a normalization map of $`𝒱`$. A reduced analytic curve is normal if and only if it is smooth. If a germ $`(𝒮,0)`$ of reduced surface is normal, then there exists a representative of it, which we keep calling $`𝒮`$, such that $`𝒮0`$ is smooth. The converse is not true. Let $`(𝒮,0)`$ be a germ of normal complex analytic surface. We say also that $`(𝒮,0)`$ is a normal surface singularity (even if the point 0 is regular on $`𝒮`$). In the sequel, we use the same notation $`(𝒮,0)`$ for the germ and for a sufficiently small representative of it. If $`e:(𝒮,0)(𝐂^N,0)`$ is any local embedding, denote by $`𝒮_{e,r}`$ the intersection of $`𝒮`$ with a euclidean ball of $`𝐂^N`$ of radius $`r1`$ and by $`M_{e,r}(𝒮)`$ the boundary of $`𝒮_{e,r}`$. By general transversality theorems due to Whitney, when $`r>0`$ is small enough, $`M_{e,r}(𝒮)`$ is a smooth manifold, naturally oriented as the boundary of the complex manifold $`𝒮_{e,r}`$. It does not depend on the choices of embedding $`e`$ and radius $`r1`$ made to define it (see Durfee ). ###### Definition 8.1. An oriented 3-manifold $`M(𝒮)`$ orientation-preserving diffeomorphic with the manifolds $`M_{e,r}(𝒮)`$, where $`r>0`$ is small enough, is called the (abstract) boundary or the link of the singularity $`(𝒮,0)`$. It is important to keep in mind that in the sequel $`M(𝒮)`$ is supposed naturally oriented as explained before. In order to understand better this remark, look at Theorem 8.11. The easiest way to describe the topological type of the manifold $`M(𝒮)`$ is (as first done by Mumford ) by retracting it to the exceptional divisor of a resolution of $`(𝒮,0)`$. Let us first define this last notion. ###### Definition 8.2. An analytic map $`p:(,E)(𝒮,0)`$ is called a resolution of the singularity $`(𝒮,0)`$ with exceptional divisor $`E=p^1(0)`$ if the following conditions are simultaneously satisfied: $``$ $``$ is a smooth surface; $``$ $`p`$ is a proper morphism; $``$ the restriction of $`p`$ to $`E=f^1(0)`$ is an isomorphism onto $`𝒮0`$. We say that $`p:(,E)(𝒮,0)`$ is a normal crossings resolution if one has moreover: $``$ $`E`$ is a divisor with normal crossings. Recall that, by definition, a divisor on a smooth complex surface has normal crossings if in the neighborhood of any of its points, its support is either smooth, or the union of transverse smooth curves. Normal crossings resolutions always exist (see Laufer and Lipman for a careful presentation of the Hirzebruch-Jung method of resolution, as well as Cossart for Zariski’s method of resolution by normalized blow-ups). There is a unique minimal resolution, which we denote $`p_{min}:(_{min},E_{min})(𝒮,0).`$ The minimality property means that any other resolution $`p:(,E)(𝒮,0)`$ can be factorized as $`p=p_{min}q`$, where $`q:_{min}`$ is a proper bimeromorphic map. The minimal resolution $`p_{min}`$ is characterized by the fact that $`E_{min}`$ contains no component $`E_i`$ which is smooth, rational and of self-intersection $`1`$ (classically called an exceptional curve of the first kind). Analogously, there is a unique resolution which is minimal among normal crossings ones. We denote it: $$p_{mnc}:(_{mnc},E_{mnc})(𝒮,0)$$ It is characterized by the fact that $`E_{mnc}`$ has normal crossings and each component $`E_i`$ of $`E_{mnc}`$ which is an exceptional curve of the first kind contains at least 3 points which are singular on $`E_{mnc}`$. If a normal crossings resolution has moreover only smooth components, one says usually that the resolution is good; there exists also a unique minimal good resolution, but in this paper we don’t consider it. The following criterion allows one to recognize the divisors which are exceptional with respect to some resolution of a normal surface singularity. ###### Theorem 8.3. Let $`E`$ be a reduced compact connected divisor in a smooth surface $``$. Denote by $`(E_i)_{1in}`$ its components. Then $`E`$ is the exceptional divisor of a resolution of a normal surface singularity if and only if the intersection matrix $`(E_iE_j)_{i,j}`$ is negative definite. The necessity is classical (see \[37, section 9\], where is presented Mumford’s proof of and where the oldest reference is to Du Val ). The sufficiency was proved by Grauert (see also Laufer ). If $`E`$ verifies the conditions which are stated to be equivalent in the theorem, one also says that $`E`$ can be contracted on $``$. From now on we suppose that $`p:(,E)(𝒮,0)`$ is a normal crossings resolution of $`(𝒮,0)`$. Denote by $`\mathrm{\Gamma }(p)`$ its weighted dual graph. Its set of vertices $`𝒱(p)`$ is in bijection with the irreducible components of $`E`$. Depending on the context, we think about $`E_i`$ as a curve on $``$ or a vertex of $`\mathrm{\Gamma }(p)`$. The vertices which represent the components $`E_i`$ and $`E_j`$ are joined by as many edges as $`E_i`$ and $`E_j`$ have intersection points on $``$. In particular, there are as many loops based at the vertex $`E_i`$ as singular points (that is, self-intersections) on the curve $`E_i`$ (see Figure 17). Each vertex $`E_i`$ is decorated by two weights, the geometric genus $`g_i`$ of the curve $`E_i`$ (that is, the genus of its normalization) and its self-intersection number $`e_i1`$ in $``$. Denote also by $`\delta _i`$ the valency of the vertex $`E_i`$, that is, the number of edges starting from it (where each loop counts for 2). For example, in Figure 17 one has $`\delta _1=9,\delta _2=5,`$ etc. ### 8.2. The plumbing structure associated to a normal crossings resolution By Definition 8.1, $`M(𝒮)`$ is diffeomorphic to $`M_{e,r}(𝒮)`$, where $`e:(𝒮,0)(𝐂^N,0)`$ is an embedding and $`r1`$. But $`M_{e,r}(𝒮)`$ is the level-set at level $`r`$ of the function $`\rho _e:(𝒮,0)(𝐑_+,0)`$, the restriction to $`e(𝒮)`$ of the distance-function to the origin in $`𝐂^N`$. As the resolution $`p`$ realizes by definition an isomorphism between $`E`$ and $`𝒮0`$, it means that $`M_{e,r}(𝒮)=\rho _e^1(r)`$ is diffeomorphic to $`\psi _e^1(r)`$, where $`\psi _e:=\rho _ep`$. The advantage of this changed viewpoint on $`M(𝒮)`$ is that it appears now orientation-preserving diffeomorphic to the boundary of a “tubular neighborhood” of the curve $`E`$ in the smooth manifold $``$. As in general $`E`$ has singularities, one has to discuss the precise meaning of the notion of tubular neighborhood. We quote Mumford \[54, pages 230-231\]: > Now the general problem, given a complex $`KE^n`$, Euclidean $`n`$-space, to define a tubular neighborhood, has been attacked by topologists in several ways although it does not appear to have been treated definitively as yet. J.H.C. Whitehead , when $`K`$ is a subcomplex in a triangulation of $`E^n`$, has defined it as the boundary of the star of $`K`$ in the second barycentric subdivision of the given triangulation. I am informed that Thom has considered it more from our point of view: for a suitably restricted class of positive $`C^{\mathrm{}}`$ fcns. $`f`$ such that $`f(P)=0`$ if and only if $`PK`$, define the tubular neighborhood of $`K`$ to be the level manifolds $`f=ϵ`$, small $`ϵ`$. The catch is how to suitably restrict $`f`$; here the archtype for $`f^1`$ may be thought of as the potential distribution due to a uniform charge on $`K`$. Let us come back to the normal crossings divisor $`E`$ in the smooth surface $``$. If $`E`$ is smooth, then one can construct a diffeomorphism between a tubular neighborhood $`U(E)`$ of $`E`$ in $``$ and of $`E`$ in the total space $`N_{}E`$ of its normal bundle in $``$. As $`N_{}E`$ is naturally fibred by discs, this is also true for $`U(E)`$. The fibration of $`U(E)`$ can be chosen in such a way that the levels $`\psi _e^1(r)`$ are transversal to the fibres for $`r1`$. In this way one gets a Seifert structure without singular fibres on $`\psi _e^1(r)M(𝒮)`$. Suppose now that $`E`$ is not smooth, but that its irreducible components are so. One can also define in this situation a notion of tubular neighborhood $`U(E)`$ of $`E`$ in $``$. One way to do it is to take the union of conveniently chosen tubular neighborhoods $`U(E_i)`$ of $`E`$’s components $`E_i`$. Abstractly, one has to glue the 4-manifolds with boundary $`U(E_i)`$ by identifying well-chosen neighborhoods of the points which get identified on $`E`$. This procedure is what is called the “plumbing” of disc-bundles over surfaces (see Hirzebruch , Hirzebruch & Neumann & Koh , Brieskorn ). Its effect on the boundaries $`U(E_i)`$ is to take out saturated filled tori and to identify their boundaries, by a diffeomorphism which permutes fibres and meridians in an orientation-preserving way. This is the 3-dimensional “plumbing” operation introduced by Mumford , alluded to in section 7.4. In order to understand what happens near a singular point of $`E`$, it is convenient to choose local coordinates $`(x,y)`$ on $`E`$ in the neighborhood of the singular point, such that $`E`$ is defined by the equation $`xy=0`$. So, $`y=0`$ defines locally an irreducible component $`E_i`$ of $`E`$ and similarly $`x=0`$ defines $`E_j`$. It is possible that $`E_i=E_j`$, a situation excluded in the previous paragraph for pedagogical reasons. If this equality is true, then the same plumbing procedure can be applied, this time by identifying well-chosen neighborhoods of points of the same 4-manifold with boundary $`U(E_i)`$. At this point appears a subtlety: the 4-manifold $`U(E_i)`$ to be considered is no longer a tubular neighborhood of $`E_i`$ in $``$, but instead of the normalization $`\stackrel{~}{E}_i`$ of $`E_i`$ inside the modified normal bundle $`\nu _i^{}T/T\stackrel{~}{E}_i`$. Here $`\nu _i:\stackrel{~}{E}_i`$ denotes the normalization map of $`E_i`$ and $`T`$, respectively $`T\stackrel{~}{E}_i`$ denote the holomorphic tangent bundles to the smooth complex manifolds $``$ and $`\stackrel{~}{E}_i`$. As a real differentiable bundle of rank 2, this vector bundle over $`\stackrel{~}{E}_i`$ is characterized by its Euler number $`\stackrel{~}{e}_i`$, which is equal to the self-intersection number of $`\stackrel{~}{E}_i`$ inside the total space of the bundle. This number is related to the self-intersection of $`E_i`$ inside $``$ in the following way (see Neumann \[56, page 333\]): ###### Lemma 8.4. If $`\stackrel{~}{e}_i`$ is the Euler number of the real bundle $`\nu _i^{}T/T\stackrel{~}{E}_i`$ over $`\stackrel{~}{E}_i`$, where $`\nu _i:\stackrel{~}{E}_i`$ is the normalization map of $`E_i`$, then $`\stackrel{~}{e}_i=e_i\delta _i`$. Proof: In order to understand this formula, just think at the effect of a small isotopy of $`E_i`$ inside $``$. Near each self-crossing point of $`E_i`$, the intersection point of one branch of $`E_i`$ with the image of the other branch after the isotopy is not counted when one computes $`\stackrel{~}{e}_i`$. $`\mathrm{}`$ Notice that Theorem 8.3 is true if one takes as diagonal entries of the matrix the numbers $`e_i=E_i^2`$, but is false if one takes instead the numbers $`\stackrel{~}{e}_i`$. The easiest example is given by an irreducible divisor $`E=E_1`$, with $`e_1=1>0`$ and $`\delta _1=2`$ which, by Lemma 8.4 implies that $`\stackrel{~}{e}_1=1<0`$. In Figure 18 we represent in two ways the local situation near the chosen singular point of $`E`$. On the left we simply draw the union of the two neighborhoods $`U(E_i)`$ and $`U(E_j)`$. On the right, “the corners are smoothed”. This is precisely what happens when we look at the levels of the function $`\psi _e`$. Moreover, we represent by interrupted lines the real analytic set defined by the equation $`\|x\|=\|y\|`$. Its intersection with $`\psi _e^1(r)U(E)M(𝒮)`$ is a two-dimensional torus $`T`$. This is the way in which such tori appear naturally as structural elements of the 3-manifolds $`M(𝒮)`$. One also sees how the complement of $`T`$ in $`U(E)`$ is fibred by boundaries of discs transversal to $`E_i`$ or $`E_j`$. By considering model neighborhoods of the singular points of $`E`$ structured as in the right-hand side of Figure 18 and conveniently extending them to a tubular neighborhood of all of $`E`$, one gets a retraction $$\mathrm{\Phi }:U(E)E$$ which restricts to a locally trivial disc-fibration over the smooth locus of $`E`$ and whose fibre over each singular point of $`E`$ is a cone over a 2-dimensional torus. By considering the restriction $`\mathrm{\Phi }\|_{U(E)}`$, we see that the fibres over the singular points of $`E`$ are embedded tori, and that their complement gets fibred by circles. As $`U(E)`$ is orientation-preserving diffeomorphic to $`M(𝒮)`$, we see that $`M(𝒮)`$ gets endowed with a graph structure $`(𝒯(p),(p))`$ well-defined up to isotopy. It is a good test of the understanding of the complexifications of Figure 18 to show that $`(𝒯(p),(p))`$ is in fact a plumbing structure (see Definition 7.15). The pieces of $`M(𝒮)_{𝒯(p)}`$ correspond to the irreducible components of $`E`$, that is to the vertices of $`\mathrm{\Gamma }(p)`$. Denote by $`M(E_i)`$ the piece which corresponds to $`E_i`$. The fibres of $`M(E_i)`$ are obtained up to isotopy by cutting the boundary of the chosen sufficiently small tubular neighborhood of $`E`$ with smooth holomorphic curves transversal to $`E`$ at smooth points of $`E_i`$. So, the plumbing structure $`(𝒯(p),(p))`$ is naturally oriented. ###### Lemma 8.5. With their natural orientations, the fibres on both sides of any component of $`𝒯(p)`$ are oriented compatibly with the orientation of $`M(𝒮)`$. Proof: The notion of compatibility we speak about is the one of Definition 7.16. We mean that, if we take an arbitrary component $`T`$ of $`𝒯(p)`$, and a tubular neighborhood $`N(T)`$ such that its preimage in $`M(𝒮)_{𝒯(p)}`$ is saturated by the leaves of the foliation $`(p)`$, then two fibres, one in each boundary component of $`N(T)`$, are oriented compatibly with the orientation of $`N(T)`$. Now, this is an instructive exercise on the geometrical understanding of the relations between the orientations of various objects in the neighborhood of a normal crossing on a smooth surface. Just think of the complexification of Figure 18. $`\mathrm{}`$ ###### Corollary 8.6. The orientation of the fibres of $`(𝒯(p),(p))`$ is determined by the associated unoriented plumbing structure up to a simultaneous change of orientation of all the fibres. Proof: Consider the unoriented plumbing structure. Start from an arbitrary piece $`M(E_i)`$, and choose one of the two continuous orientations of its fibres. Then propagate this orientations farther and farther through the components of $`𝒯(p)`$, by respecting the compatibility condition on the neighboring orientations. As $`M(𝒮)`$ is connected, we know that after a finite number of steps one has oriented the fibres of all the pieces. As one orientation exists which is compatible in the neighborhood of all the tori, we see that our process cannot arrive at a contradiction (that is, a non-trivial monodromy around a loop of $`\mathrm{\Gamma }(p)`$ in the choice of orientations). $`\mathrm{}`$ The following lemma is a particular case of the study done in Mumford \[54, page 11\] and Hirzebruch \[35, page 250-03\]. ###### Lemma 8.7. Suppose that $`E_i`$ is a component of $`E`$ which is smooth, rational and whose valency in the graph $`\mathrm{\Gamma }(p)`$ is $`2`$. In the thick torus $`M(E_i)`$ which corresponds to it in the plumbing structure $`(𝒯(p),(p))`$, consider an oriented fibre $`f`$ of $`M(E_i)`$, as well as oriented fibres $`f^{}`$, $`f^{\prime \prime }`$ of the two (possibly coinciding) adjacent pieces. Then one has the following relation in the homology group $`H_1(M(E_i),𝐙)`$: $$[f^{}]+[f^{\prime \prime }]=\|e_i\|[f].$$ ### 8.3. The topological characterization of HJ and cusp singularities We want now to understand how to pass from the plumbing structure $`(𝒯(p),(p))`$ on $`M(𝒮)`$ to the canonical graph structure on it (see Definition 7.14). We see that the pieces of $`M(𝒮)_{𝒯(p)}`$ which are thick tori correspond to components $`E_i`$ which are smooth and rational with $`\delta _i=2`$, and those which are solid tori correspond to components $`E_i`$ which are smooth and rational with $`\delta _i=1`$. It is then natural to introduce the following: ###### Definition 8.8. We say that a vertex $`E_i`$ of $`\mathrm{\Gamma }(p)`$ is a chain vertex if $`E_i`$ is smooth, $`g_i=0`$ and $`\delta _i2`$. If moreover $`\delta _i=2`$, we call it an interior chain vertex, otherwise we call it a terminal chain vertex. We say that a vertex of $`\mathrm{\Gamma }(p)`$ is a node if it is not a chain vertex. In , Lê, Michel & Weber used the name “rupture vertex” for a node in the dual graph associated to the minimal embedded resolution of a plane curve singularity. In their situation, where all the vertices represent smooth rational curves, nodes are simply those of valency $`3`$. In our case this is no longer true, as one can have also vertices of valency $`2`$, if they correspond to curves $`E_i`$ which are either not smooth or of genus $`g_i1`$. Denote by $`𝒩(p)`$ the set of nodes of $`\mathrm{\Gamma }(p)`$. It is an empty set if and only if $`\mathrm{\Gamma }(p)`$ is topologically a segment or a circle and all the components $`E_i`$ are smooth rational curves. The first situation occurs precisely for the Hirzebruch-Jung singularities, defined in Section 6.2 (see Proposition 6.2), and the second one for cusp singularities, introduced by Hirzebruch in the number-theoretical context of the study of Hilbert modular surfaces. ###### Definition 8.9. A germ $`(𝒮,0)`$ of normal surface singularity is called a cusp singularity if it has a resolution $`p`$ such that $`\mathrm{\Gamma }(p)`$ is topologically a circle and $`𝒩(p)=\mathrm{}`$. For other definitions and details about them, see Hirzebruch , Laufer (where they appear as special cases of minimally elliptic singularities), Ebeling & Wall (where they appear as special cases of Kodaira singularities), Oda , Wall and Némethi . They were generalized to higher dimensions by Tsuchihashi (see Oda \[59, Chapter 4\]). In the previous definition it is not possible to replace the resolution $`p`$ by the minimal normal crossings one. Indeed: ###### Lemma 8.10. If $`(𝒮,0)`$ is a cusp singularity, then $`\mathrm{\Gamma }(p_{mnc})`$ is topologically a circle and either $`𝒩(p_{mnc})=\mathrm{}`$, or $`E_{mnc}`$ is irreducible, rational, with one singular point where it has normal crossings. Proof: One passes from $`p`$ to $`p_{mnc}`$ by successively contracting components $`F`$ which are smooth, rational and verify $`F^2=1`$ (that is, exceptional curves of the first kind, by a remark which follows Definition 8.2). The new exceptional divisor verifies the same hypothesis as the one of $`p`$, except when one passes from a divisor with 2 components to a divisor with one component. In this last situation, this second irreducible divisor is rational, as its strict transform $`F`$ is so. Moreover, it has one singular point with normal crossing branches passing through it, as by hypothesis $`F`$ cuts transversely the other component of the first divisor in exactly two points. $`\mathrm{}`$ We would like to emphasize the following theorem due to Neumann \[56, Theorem 3\], which characterizes Hirzebruch-Jung and cusp singularities among normal surface singularities. ###### Theorem 8.11. Let $`(𝒮,0)`$ be a normal surface singularity. The manifold $`M(𝒮)`$ is orientation-preserving diffeomorphic to the abstract boundary of a normal surface singularity if and only if $`(𝒮,0)`$ is either a Hirzebruch-Jung singularity or a cusp-singularity. Recall that $`M(𝒮)`$ denotes the manifold $`M(𝒮)`$ with reversed orientation. We will bring more light on this theorem with Propositions 9.3 and 9.6, which show that for both Hirzebruch-Jung and cusp singularities, the involutions $`M(𝒮)M(𝒮)`$ are manifestations of the duality described in section 5. As Hirzebruch-Jung singularities, cusp singularities can also be defined using toric geometry (see Oda \[59, Chapter 4\]). In the same spirit, as a particular case of Laufer’s classification of taut singularities, we have: ###### Theorem 8.12. Hirzebruch-Jung and cusp singularities are taut, that is, their analytical type is determined by their topological type. For this reason, it is natural to ask which 3-manifolds are obtained as abstract boundaries of Hirzebruch-Jung singularities and cusp singularities. This question is answered by: ###### Proposition 8.13. 1) $`(𝒮,0)`$ is a Hirzebruch-Jung singularity if and only if $`M(𝒮)`$ is a lens space. Moreover, each oriented lens space appears like this. 2) $`(𝒮,0)`$ is a cusp singularity if and only if $`M(𝒮)`$ is a torus fibration with algebraic monodromy of trace $`3`$. Moreover, each oriented torus fibration of this type appears like this. Proof: This proposition is a particular case of Neumann \[56, Corollary 8.3\]. Here we sketch the proofs of the necessities, in order to develop tools for sections 9.1 and 9.2. Let $`p:(,E)(𝒮,0)`$ be the minimal normal crossings resolution of $`(𝒮,0)`$ (for notational convenience, we drop the index “mnc”). Denote by $`U(E)`$ a (closed) tubular neighborhood of $`E`$ in $``$ and by $`\mathrm{\Phi }:U(E)E`$ a preferred retraction, as defined in section 8.2. Denote also by $$\mathrm{\Psi }:U(E)E$$ the restriction of $`\mathrm{\Phi }`$ to $`U(E)M(𝒮)`$. 1) Suppose that $`(𝒮,0)`$ is a Hirzebruch-Jung singularity. Orient the segment $`\mathrm{\Gamma }(p)`$. Denote then by $`E_1,\mathrm{},E_r`$ the components of $`E`$ in the order in which they appear along $`\mathrm{\Gamma }(p)`$ in the positive direction. For each $`i\{1,\mathrm{},r1\}`$, denote by $`A_{i,i+1}`$ the intersection point of $`E_i`$ and $`E_{i+1}`$. Consider also two other points $`A_{0,1}E_1,A_{r,r+1}E_r`$ which are smooth points of $`E`$. Then consider on each component $`E_i`$ a Morse function $$\mathrm{\Pi }_i:E_i[\frac{i1}{r},\frac{i}{r}]$$ having as its only critical points $`A_{i1,i}`$ (where $`\mathrm{\Pi }_i`$ attains its minimum) and $`A_{i,i+1}`$ (where $`\mathrm{\Pi }_i`$ attains its maximum). As $`\mathrm{\Pi }_i(A_{i,i+1})=\mathrm{\Pi }_{i+1}(A_{i,i+1})`$ for all $`i\{1,\mathrm{},r1\}`$, we see that the maps $`\mathrm{\Pi }_i`$ can be glued together in a continuous map $$\mathrm{\Pi }:E[0,1].$$ Consider the composed continuous map $`\mathrm{\Pi }\mathrm{\Psi }:M(𝒮)[0,1]`$ (see Figure 19). Our construction shows that its fibres over $`0`$ and $`1`$ are circles and that those over interior points of $`[0,1]`$ are tori. Moreover, each such torus splits $`M`$ into two solid tori. By Definition 7.3, we see that $`M`$ is a lens space. It remains now to prove that each oriented lens space appears like this. Denote $`L:=H_1(M(𝒮)(\mathrm{\Pi }\mathrm{\Psi })^1\{0,1\},𝐙)`$. As $`M(𝒮)(\mathrm{\Pi }\mathrm{\Psi })^1\{0,1\}`$ is the interior of a thick torus foliated by the tori $`(\mathrm{\Pi }\mathrm{\Psi })^1(c)`$, where $`c(0,1)`$, we see that $`L`$ is a 2-dimensional lattice. With the notations of section 8.2, let $`f_i`$ be an oriented fibre in the piece $`M(E_i)`$ of the plumbing structure $`(𝒯(p),(p))`$ which corresponds to $`E_i`$. Consider also $`f_0`$ and $`f_{r+1}`$, canonically oriented meridians on the boundaries of tubular neighborhoods of $`(\mathrm{\Pi }\mathrm{\Psi })^1(0)`$, respectively $`(\mathrm{\Pi }\mathrm{\Psi })^1(1)`$. For each $`i\{0,\mathrm{},r+1\}`$, denote by $`v_i:=[f_i]L`$ the homology class of $`f_i`$. Recall that $`e_i:=E_i^2`$. By Lemma 8.7, we see that (25) $$v_{i+1}=e_iv_iv_{i1},i\{0,\mathrm{},r\}.$$ By Proposition 6.2, $`p`$ is also the minimal resolution of $`(𝒮,0)`$, which shows that $`\|e_i\|2,i\{1,\mathrm{},r\}`$. Now apply Proposition 4.4. We deduce that the numbers $`e_i`$ are determined by the oriented topological type of the lens space $`M(𝒮)`$, once the isotopy class of the tori $`(\mathrm{\Pi }\mathrm{\Psi })^1(c)`$ is fixed. This shows that, starting from any oriented lens space $`M`$ and torus $`TM`$ which splits $`M`$ into two solid tori, one can construct a Hirzebruch-Jung singularity $`(𝒮,0)`$ such that $`M(𝒮)M`$ only by looking at the classes of the meridians of the two solid tori in the lattice $`L=H_1(T,𝐙)`$. One has only to be careful to orient them compatibly with the orientation of $`M`$ (as explained at the beginning of the proof of Lemma 8.5). 2) Suppose that $`(𝒮,0)`$ is a cusp singularity. $``$ Consider first the case where $`r2`$. Orient the circle $`\mathrm{\Gamma }(p)`$ and choose one of its vertices. Denote then by $`E_1,\mathrm{},E_r`$ the components of $`E`$ in the order in which they appear along $`\mathrm{\Gamma }(p)`$ in the positive direction, starting from $`E_1`$. For each $`i\{1,\mathrm{},r\}`$, denote by $`A_{i,i+1}`$ the intersection point of $`E_i`$ and $`E_{i+1}`$, where $`E_{r+1}=E_1`$. Consider then functions $`\mathrm{\Pi }_i:E_i[\frac{i1}{r},\frac{i}{r}]`$ with the same properties as in the case of Hirzebruch-Jung singularities. By passing to the quotient $`𝐑𝐑/𝐙`$, we can glue the previous maps into a continuous map: $$\mathrm{\Pi }:E𝐑/𝐙.$$ Consider then the map $`\mathrm{\Pi }\mathrm{\Psi }:M(𝒮)𝐑/𝐙`$ (see Figure 20). Our construction shows that $`\mathrm{\Pi }`$ realizes $`M(𝒮)`$ as the total space of a torus fibration over $`𝐑/𝐙`$. Denote by $`T_{i,i+1}:=\mathrm{\Psi }^1(A_{i,i+1})`$ the torus of $`𝒯(p)`$ which corresponds to the intersection point of $`E_i`$ and $`E_{i+1}`$. Denote $`T:=T_{r,1}`$ and let $`N(T)`$ be a (closed) tubular neighborhood of $`T`$, which does not intersect any other torus $`T_{i,i+1}`$, for $`i\{1,\mathrm{},r1\}`$ (see Figure 20). Denote $`L:=H_1(M(𝒮)N(T),𝐙)`$. As $`M(𝒮)N(T)`$ is the interior of a thick torus, we see that $`L`$ is a 2-dimensional lattice. With the notations of section 8.2, let $`f_i`$ be an oriented fibre in the piece $`M(E_i)`$. We suppose moreover that $`f_1`$ and $`f_r`$ are situated on the boundary of $`N(T)`$. Consider two other circles $`f_0`$ and $`f_{r+1}`$ on $`N(T)`$, such that $`f_0,f_r`$ are isotopic inside $`N(T)`$ and situated on distinct boundary components and such that the same is true for the pair $`f_1,f_{r+1}`$. For each $`i\{0,\mathrm{},r+1\}`$, denote by $`v_i:=[f_i]L`$ the homology class of $`f_i`$. By Lemma 8.7, we see that: (26) $$v_{i+1}=e_iv_iv_{i1}=e_iv_iv_{i1},i\{0,\mathrm{},r\},$$ where $`E_0:=E_r`$. Denote by $`nGL(L)`$ the automorphism which sends the basis $`(v_0,v_1)`$ of $`L`$ into the basis $`(v_r,v_{r+1})`$. The relations (26) show that its matrix in the basis $`(v_0,v_1)`$ is: $$\left(\begin{array}{cc}0& 1\\ 1& e_1\end{array}\right)\left(\begin{array}{cc}0& 1\\ 1& e_2\end{array}\right)\mathrm{}\left(\begin{array}{cc}0& 1\\ 1& e_r\end{array}\right)$$ A little thinking shows that $`n`$ is the inverse of the algebraic monodromy $`mGL(L)`$ in the positive direction along $`𝐑/𝐙`$. So, the matrix of $`m`$ in the basis $`(v_0,v_1)`$ is: $$\left(\begin{array}{cc}e_r& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}e_{r1}& 1\\ 1& 0\end{array}\right)\mathrm{}\left(\begin{array}{cc}e_1& 1\\ 1& 0\end{array}\right)$$ We have reproved like this Theorem 6.1 IV in Neumann . We deduce by induction the following expression for its trace, where the polynomials $`Z^{}`$ were defined by formula (1): (27) $$trm=Z^{}(\|e_1\|,\mathrm{},\|e_r\|)Z^{}(\|e_2\|,\mathrm{},\|e_{r1}\|).$$ The negative definiteness of the intersection matrix of $`E`$ (see Theorem 8.3) shows that there exists $`i\{1,\mathrm{},r\}`$ such that $`\|e_i\|3`$. As $`p`$ is supposed to be the minimal resolution of $`(𝒮,0)`$, we have also $`e_j2,j\{1,\mathrm{},r\}`$. Using equation (27), we deduce then easily by induction on $`r`$ that $`trm3`$. $``$ Consider now the case $`r=1`$. Then, by Lemma 8.10, $`E`$ is a rational curve with one singular point $`P`$, where $`E`$ has normal crossings. Let $`p^{}:(^{},E^{})(𝒮,0)`$ be the resolution of $`(𝒮,0)`$ obtained by blowing up $`P`$. Then $`E^{}`$ is a normal crossings resolution with smooth components $`E_1,E_2`$, where $`E_1^2=1`$ and $`E_2`$ is the strict transform of $`E`$. As $`(p^{})^{}E=2E_1+E_2`$ and $`((p^{})^{}E)^2=E^2`$, we deduce that $`E_2^2=E^245`$. Now we apply the same argument as in the case $`r2`$, but for the resolution $`p^{}`$. An alternative proof could use Lemma 8.4. The fact that each oriented torus fibration with $`trm3`$ appears like this is a consequence of the study done in section 9.2. Indeed, there we show how to extract the numbers $`(e_1,\mathrm{},e_r)`$ from the oriented topological type of $`M(𝒮)`$. $`\mathrm{}`$ By Neumann , there exist also abstract boundaries $`M(𝒮)`$ which are torus fibrations with algebraic monodromy of trace $`2`$. But in that case the exceptional divisor of the minimal resolution is an elliptic curve (then, following Saito , one speaks about simple elliptic singularities, which are other particular cases of minimally elliptic ones). ### 8.4. Construction of the canonical graph structure Consider again an arbitrary normal surface singularity $`(𝒮,0)`$ and a normal crossings resolution $`p`$ of it. ###### Definition 8.14. Suppose that the set of nodes $`𝒩(p)`$ is non-empty. Conceive the graph $`\mathrm{\Gamma }(p)`$ as a 1-dimensional CW-complex and take the complement $`\mathrm{\Gamma }(p)𝒩(p)`$. This complement is the disjoint union of segments, which we call chains. If a chain is open at both extremities we call it an interior chain. If it is half-open we call it a terminal chain. In Figure 20 we represent the chains of Figure 17, with the hypothesis that $`E_4,E_5,E_7𝒩(p)`$ and $`E_6𝒩(p)`$. That is, we suppose that $`E_4,E_5,E_6,E_7`$ are smooth and that $`g(E_4)=g(E_5)=g(E_7)=0,g(E_6)1`$. There is only one terminal chain, which contains the terminal chain vertex $`E_7`$. Denote by $`𝒞(p)`$ the set of chains. This set can be written as a disjoint union $$𝒞(p)=𝒞_i(p)𝒞_t(p)$$ where $`𝒞_i(p)`$ denotes the set of interior chains and $`𝒞_t(p)`$ the set of terminal chains. The edges of $`\mathrm{\Gamma }(p)`$ contained in a chain $`C𝒞(p)`$ correspond to a set of parallel tori in $`M(𝒮)`$. Choose one torus $`T_C`$ among them and define: $$𝒯^{}(p):=\underset{C𝒞_i(p)}{}T_C.$$ By construction, each piece of $`M(𝒮)_{𝒯^{}(p)}`$ contains a unique piece $`M(E_i)`$ of $`M(𝒮)_{𝒯(p)}`$ such that $`E_i`$ is a node of $`\mathrm{\Gamma }(p)`$. If $`E_i𝒩(p)`$, denote by $`M^{}(E_i)`$ the piece of $`M(𝒮)_{𝒯^{}(p)}`$ which contains $`M(E_i)`$. One can extend in a unique way up to isotopy the natural Seifert structure without exceptional fibres on $`M(E_i)`$ to a Seifert structure on $`M^{}(E_i)`$. One obtains like this a graph structure $`(𝒯^{}(p),^{}(p))`$ on $`M(𝒮)`$. Till now we have worked with any normal crossings resolution $`p`$. We consider now a special one, the minimal normal crossings resolution $`p_{mnc}`$. ###### Proposition 8.15. Suppose that $`(𝒮,0)`$ is neither a Hirzebruch-Jung singularity, nor a cusp singularity. Then the graph structure $`(𝒯^{}(p_{mnc}),^{}(p_{mnc}))`$ is the canonical graph structure on $`M(𝒮)`$. Proof: If $`𝒯^{}(p_{mnc})`$ is empty, as $`(𝒮,0)`$ is not a cusp singularity we deduce that $`(𝒯^{}(p_{mnc}),^{}(p_{mnc}))`$ is a Seifert structure. By Proposition 7.8, we see that it is the canonical graph structure on $`M(𝒮)`$. Suppose now that $`𝒯^{}(p_{mnc})`$ is non-empty. One has to verify two facts (see Definition 7.14): $``$ first, that all the fibrations induced by $`^{}(p_{mnc})`$ on the pieces which are thick Klein bottles have orientable basis; $``$ second, that by taking the various choices of Seifert structures on the pieces of $`M(𝒮)_{𝒯^{}(p_{mnc})}`$, one does not obtain isotopic fibres coming from different sides on one of the tori of $`𝒯^{}(p_{mnc})`$. The first fact is immediate, as one starts from Seifert structures with orientable basis on the pieces of $`M(𝒮)_{𝒯(p_{mnc})}`$ before eliminating tori of $`𝒯(p_{mnc})`$ in order to remain with $`𝒯^{}(p_{mnc})`$. In what concerns the second fact, the idea is to look at the fibres corresponding to the chain vertices of any interior chain $`C`$. The union of the pieces of $`M(𝒮)_{𝒯(p_{mnc})}`$ which are associated to those vertices is a thick torus $`N_R`$. Take a fibre in each piece (remember that they are naturally oriented as boundaries of holomorphic discs) and look at their images in $`L=H_1(N_R,𝐙)`$. One gets like this a sequence of vectors $`v_1,\mathrm{},v_sL`$. Consider also the images $`v_0`$ and $`v_{s+1}`$ of the fibres coming from the nodes of $`\mathrm{\Gamma }(p_{mnc})`$ to which $`C`$ is adjacent, the order of the indices respecting the order of the vertices along the chain. By Lemma 8.7, $`v_{k+1}=\alpha _kv_kv_{k1}`$ for any $`k\{1,\mathrm{},s\}`$, where $`\alpha _k`$ is the absolute value of the self-intersection of the component $`E_i`$ of $`\mathrm{\Gamma }(p_{mnc})`$ which gave rise to the vector $`v_k`$. Here plays the hypothesis that $`p_{mnc}`$ is minimal: this implies that $`\alpha _k2`$. Then one can conclude by using Proposition 7.7. The analysis of thick Klein bottles is similar. It is based on the fact that a thick Klein bottle can appear only from a portion of the graph $`\mathrm{\Gamma }(p)`$ as in Figure 21, where $`E_1,E_2,E_3`$ are smooth rational curves of self-intersections $`2,2`$, respectively $`n`$ (see Neumann \[56, pages 305, 334\]). The important point is that $`n2`$. Otherwise the complete sub-graph of $`\mathrm{\Gamma }(p)`$ with vertices $`E_1,E_2,E_3`$ would have a non-definite intersection matrix, which contradicts Theorem 8.3. $`\mathrm{}`$ The plumbing structure $`(𝒯(p_{mnc}),(p_{mnc}))`$ on $`N(𝒮)`$ is associated to the resolution $`p_{mnc}`$ of $`(𝒮,0)`$. One can wonder if the canonical graph structure $`(𝒯^{}(p_{mnc}),`$ $`^{}(p_{mnc}))`$ is also associated to some analytic morphism with target $`(𝒮,0)`$. This is indeed the case. In order to see it, start from $`p_{mnc}`$ and its exceptional divisor $`E`$. Then contract all the components of $`E`$ which correspond to chain vertices. One gets like this a normal surface with only Hirzebruch-Jung singularities. The image of $`E`$ on it is a divisor $`F`$ with again only normal crossings when seen as an abstract curve. Take then as a representative of $`M(𝒮)`$ the boundary of a tubular neighborhood of $`F`$ in the new surface and split it into pieces which project into the various components of $`F`$. The splitting is done using tori which are associated bijectively to the singular points of $`F`$. Namely, in a system of (toric) local coordinates $`(x,y)`$ such that $`F`$ is defined by $`xy=0`$, one proceeds as for the definition of the plumbing structure associated to a normal crossings resolution (see Section 8.2). Then this system of tori is isotopic to $`𝒯^{}(p_{mnc})`$. ## 9. Invariance of the canonical plumbing structure on the boundary of a normal surface singularity In this section we describe how to reconstruct the plumbing structure $`(𝒯(p_{mnc}),(p_{mnc}))`$ on $`M(𝒮)`$ associated to the minimal normal crossings resolution of $`(𝒮,0)`$, only from the abstract oriented manifold $`M(𝒮)`$. Namely, using the classes of plumbing structures on thick tori defined in section 7.5, we define a plumbing structure $`𝒫(M(𝒮))`$ on $`M(𝒮)`$ and we prove: ###### Theorem 9.1. 1) When considered as an unoriented structure, the plumbing structure $`𝒫(M(𝒮))`$ depends up to isotopy only on the natural orientation of $`M(𝒮)`$. We call it the canonical plumbing structure on $`M(𝒮)`$. 2) The plumbing structure $`(𝒯(p_{mnc}),(p_{mnc}))`$ associated to the minimal normal crossings resolution of $`(𝒮,0)`$ is isotopic to the canonical plumbing structure $`𝒫(M(𝒮))`$. As a corollary we get the theorem of invariance of the plumbing structure $`(𝒯(p_{mnc}),(p_{mnc}))`$ announced in the introduction (see Theorem 9.7). We also explain how the orientation reversal on the boundary of a Hirzebruch-Jung or cusp singularity reflects the duality between supplementary cones explained in section 5.1 (see Propositions 9.4 and 9.6). In order to prove Theorem 9.1, we consider three cases, according to the nature of $`M(𝒮)`$. In the first one it is supposed to be a lens space, in the second one a torus fibration with algebraic monodromy of trace $`3`$ and in the last one none of the two (so, by Proposition 8.13, this corresponds to the trichotomy: $`(𝒮,0)`$ is a Hirzebruch-Jung singularity/ a cusp singularity/ none of the two). The idea is to start from some structure on $`M(𝒮)`$ which is well-defined up to isotopy, and to enrich it by canonical constructions of Hirzebruch-Jung plumbing structures (defined in section 7.5). When $`M(𝒮)`$ is neither a lens space nor a torus fibration with algebraic monodromy of trace $`3`$, this starting structure will be the canonical graph structure (see Definition 7.14). Otherwise we need some special theorems of structure (Theorems 9.2 and 9.5). ### 9.1. The case of lens spaces Notice that by Proposition 8.13 1), $`M(𝒮)`$ is a lens space if and only if $`(𝒮,0)`$ is a Hirzebruch-Jung singularity. The following theorem was proved by Bonahon : ###### Theorem 9.2. Up to isotopy, a lens space contains a unique torus which splits it into two solid tori. We say that a torus embedded in a lens space and splitting it into two solid tori is a central torus. By the previous theorem, a central torus is well-defined up to isotopy. Let $`M`$ be an oriented lens space and $`T`$ a central torus in $`M`$. Consider a tubular neighborhood $`N(T)`$ of $`T`$ in $`M`$, whose boundary components we denote by $`T_{}`$ and $`T_+`$, ordered in an arbitrary way. Then $`M_{T_{}T_+}`$ has three pieces, one being sent diffeomorphically by the reconstruction map $`r_{M,T_{}T_+}`$ on $`N(T)`$ \- by a slight abuse of notations, we keep calling it $`N(T)`$ \- and the others, $`M_{}`$ and $`M_+`$, having boundaries sent by $`r_{M,T_{}T_+}`$ on $`T_{}`$, respectively $`T_+`$ (see Figure 33). The manifolds $`M_{}`$ and $`M_+`$ are solid tori, as $`T`$ was supposed to be a central torus. Let $`\gamma _{}`$ and $`\gamma _+`$ be meridians of $`M_{}`$, respectively $`M_+`$, oriented compatibly with the orientation of $`N(T)`$ (see Definition 7.16). Consider the Hirzebruch-Jung plumbing structure $`𝒫(N(T),\gamma _{},\gamma _+)`$ on $`N(T)`$, whose tori are denoted by $`T_0=T_{},T_1,\mathrm{},T_{r+2}=T_+`$, as explained in Section 7.5. Denote $`𝒯_M:=T_2\mathrm{}T_r`$. Then $`M_{𝒯_M}`$ contains four pieces less than the manifold $`M_{T_{}𝒯_MT_+}`$. Denote by $`M_{}^{}`$ and $`M_+^{}`$ the piece which “contains” $`M_{}`$, respectively $`M_+`$. On $`M_{}^{}`$ we consider the Seifert structure which extends the Seifert structure of $`M_1`$ and on $`M_+^{}`$ the one which extends the Seifert structure of $`M_r`$. By applying the intersection theoretical criterion of Proposition 7.7 b), we see that those Seifert structures have no exceptional fibres (we used a similar argument to construct in Section 7.5 the Hirzebruch-Jung plumbing structure on solid tori). On the other pieces of $`M_{𝒯_M}`$ we consider the Seifert structure coming from the plumbing structure $`𝒫(M,\gamma _{},\gamma _+)`$. Denote by $`𝒫(M)`$ the plumbing structure constructed like this on the oriented manifold $`M`$. Proof of Theorem 9.1: 1) This is obvious by construction (we use Theorem 9.2). 2) In the construction of $`𝒫(M(𝒮))`$, one can take as central torus $`T`$ any torus $`(\mathrm{\Pi }\mathrm{\Psi })^1(c)`$, with $`c(0,1)`$, in the notations of the proof of Proposition 8.13, 1). Then one sees that $`[\gamma _{}]=[f_0]`$ and $`[\gamma _+]=[f_{r+1}]`$ in the lattice $`L=H_1(M(𝒮)(\mathrm{\Pi }\mathrm{\Psi })^1\{0,1\},𝐙)=H_1(T,𝐙)`$. Using the relations (25) and the definition of a Hirzebruch-Jung plumbing structure on a thick torus (see section 7.5), we deduce that the images of the fibres $`f_i`$ in $`L`$ are equal to the images of the fibres of $`𝒫(M(𝒮))`$ (see also Proposition 4.4). The proposition follows by the fact that on a 2-torus, any oriented essential curve is well-defined up to isotopy by its homology class. $`\mathrm{}`$ Let $`\sigma `$ be the strictly convex cone of $`L_𝐑`$ whose edges are generated by $`[\gamma _{}]`$ and $`[\gamma _+]`$. If one changes the ordering of the components of $`N(T)`$, then one gets the same cone $`\sigma `$, and if one changes simultaneously the orientations of $`\gamma _{}`$ and $`\gamma _+`$, then one gets the opposite cone. But if one changes the orientation of $`M`$, then the cone $`\sigma `$ is replaced by a supplementary cone. So, in view of Section 5.3, the two cones are in duality. In this sense, the canonical plumbing structure $`𝒫(M(𝒮))`$ is dual to $`𝒫(M(𝒮))`$. We get: ###### Proposition 9.3. Let $`(𝒮,0)`$ be a Hirzebruch-Jung singularity. Then the canonical plumbing structures with respect to the two possible orientations of $`M(𝒮)`$ are dual to each other. More precisely, if $`(𝒮,0)(𝒵(L,\sigma ),0)𝒜_{p,q}`$, then $`M(𝒮)`$ is orientation preserving diffeomorphic to $`M(\stackrel{ˇ}{𝒮})`$, where, with the notations of section 4, $`(\stackrel{ˇ}{𝒮},0)(𝒵(\stackrel{ˇ}{L},\stackrel{ˇ}{\sigma }),0)𝒜_{p,pq}`$. Let $`\lambda :={\displaystyle \frac{p}{q}}`$ be the type of the cone $`(L,\sigma )`$ in the sense of Definition 5.5, where $`0<q<p`$ and $`gcd(p,q)=1`$. The oriented lens space $`M(𝒮)`$, where $`(𝒮,0)(𝒵(L,\sigma ),0)𝒜_{p,q}`$, is said classically to be of type $`L(p,q)`$. By Propositions 5.6 and 5.8, combined with Theorem 9.2, we get the following classical fact: ###### Proposition 9.4. 1) The lens spaces $`L(p,q)`$ and $`L(p,q^{})`$ are orientation-preserving diffeomorphic if and only if $`p=p^{}`$ and $`q^{}\{q,\overline{q}\}`$, where $`0<\overline{q}<p,q\overline{q}1(modp)`$. 2) The lens spaces $`L(p,q)`$ and $`L(p,q^{})`$ are orientation-reversing diffeomorphic if and only if $`p=p^{}`$ and $`q^{}\{pq,p\overline{q}\}`$. ### 9.2. The case of torus fibrations with $`trm3`$ Notice that by Proposition 8.13 2), $`M(𝒮)`$ is a torus fibration whose algebraic monodromy verifies $`trm3`$ if and only if $`(𝒮,0)`$ is a cusp singularity. First we study with a little more detail torus fibrations. Let $`M`$ be an orientable torus fibration. Take a fibre torus $`T`$. Then consider the lattice $`L=H_1(T,𝐙)`$ and the algebraic monodromy operator $`mSL(L)`$ (see Definition 7.4) associated with one of the two possible orientations of the base. The following theorem is a consequence of Waldhausen \[82, section 3\] (see also Hatcher \[33, section 5\]): ###### Theorem 9.5. Up to isotopy, an orientable torus fibration $`M`$ such that $`trm3`$ contains a unique torus which splits it into a thick torus (see Definition 7.2). We say that a torus embedded in an orientable torus fibration whose algebraic monodromy $`m`$ verifies $`trm3`$ and which splits it into a thick torus is a fibre torus. By the previous theorem, a fibre torus is well-defined up to isotopy. From now on, we suppose that indeed $`trm3`$ (see Proposition 8.13, 2)). As $`M`$ is orientable, $`m`$ preserves the orientation of $`L`$, which shows that $`detm=1`$. This implies that the characteristic polynomial of $`m`$ is $`X^2(trm)X+1`$. We deduce that $`m`$ has two strictly positive eigenvalues with product $`1`$, and so the eigenspaces are two distinct real lines in $`L_𝐑`$. But the most important point is that these lines are irrational. Indeed, the eigenvalues are $`\nu _\pm :=\frac{1}{2}(trm\pm \sqrt{(trm)^24})`$ and $`(trm)^24`$ is never a square if $`trm3`$. Denote by $`d_{}`$ and $`d_+`$ the eigenspaces corresponding to $`\nu _{}`$, respectively $`\nu _+`$. Then $`m`$ is strictly contracting when restricted to $`d_{}`$ and strictly expanding when restricted to $`d_+`$. Choose arbitrarily one of the two half-lines in which $`0`$ divides the line $`d_{}`$, and call it $`l_{}`$. At this point we have not used any orientation of $`M`$. Suppose now that $`M`$ is oriented. Then the chosen orientation on the basis of the torus fibration induces an orientation of the fibre torus $`T`$, by deciding that this orientation, followed by the transversal orientation which projects on the orientation of the base induces the ambient orientation on $`M`$. Denote by $`l_+`$ the half-line bounded by $`0`$ on $`d_+`$ into which $`l_{}`$ arrives first when turned in the negative direction. Let $`\sigma `$ be the strictly convex cone bounded by these two half-lines (see Figure 24). We arrive like this at a pair $`(L,\sigma )`$ where both edges of $`\sigma `$ are irrational. As $`m`$ preserves $`L`$ and $`\sigma `$, it preserves also the polygonal line $`P(\sigma )`$. Let $`P_1`$ be an arbitrary integral point of $`P(\sigma )`$. Consider the sequence $`(P_n)_{n1}`$ of integral points of $`P(\sigma )`$ read in the positive direction along $`P(\sigma )`$, starting from $`P_1`$. There exists an index $`t1`$ such that $`P_{t+1}=m(P_1)`$. It is the period of the action of $`m`$ on the linearly ordered set of integral points of $`P(\sigma )`$. Consider $`t`$ parallel tori $`T_1,\mathrm{},T_t`$ inside $`M`$, where $`T_1=T`$ and the indices form an increasing function of the orders of appearance of the tori when one turns in the positive direction. Denote $`𝒯:=_{1kt}T_k`$ and $`T_{t+1}:=T_1`$. For each $`k\{1,\mathrm{},t\}`$, denote by $`M_k`$ the piece of $`M_𝒯`$ whose boundary components project by $`r_{M,𝒯}`$ on $`T_k`$ and $`T_{k+1}`$ (see Figure 25). Then look at the thick torus $`M_T`$. Let $`T_{}`$ be its boundary component through which one “enters inside” $`M_T`$ when one turns in the positive direction, and $`T_+`$ be the one by which one “leaves” $`M_T`$. Identify then $`H_1(M_T,𝐙)`$ with $`H_1(T_{},𝐙)`$ through the inclusion $`T_{}M_{T_1}`$ and $`H_1(T_{},𝐙)`$ with $`H_1(T,𝐙)=L`$ through the reconstruction mapping $`r_{M,T}\|_T_{}:T_{}T`$. Consider now on each piece $`M_k`$ an oriented Seifert fibration $`_k`$ such that the class of a fibre in $`L`$ (after projection in $`M_T`$ and identification of $`H_1(M_T,𝐙)`$ with $`L`$, as explained before) is equal to $`OP_k`$. Denote by $``$ the Seifert structure on $`M_𝒯`$ obtained by taking the union of the structures $`_k`$. We get like this a plumbing structure on $`M`$. Denote it by $`𝒫(M)`$. This plumbing structure does not depend, up to isotopy, on the choice of the initial integral point on $`P(\sigma )`$. Indeed, by lifting to $`M`$ a vector field of the form $`\frac{}{\theta }`$ on the base of the torus fibration and by considering its flow, one sees that one gets isotopic torus fibrations by starting from any integral point of $`P(\sigma )`$. Notice that it does neither depend on the choice of the half-line $`l_{}`$. An opposite choice would lead to the choice of an opposite cone, that is to the same unoriented plumbing structure. Proof of Theorem 9.1: 1) This is obvious by construction (we use Theorem 9.5). 2) In the construction of $`𝒫(M(𝒮))`$, one can take as fibre torus $`T`$ the torus $`T_{r,1}`$, with the notations of the proof of Proposition 8.13, 2). Using the relations (26) and Proposition 4.4, we get the Proposition. $`\mathrm{}`$ By Theorem 8.12, cusp singularities are determined up to analytic isomorphism by the topological type of the oriented manifold $`M(𝒮)`$. By Theorem 9.5, this manifold can be encoded by a pair $`(T,\mu )`$, where $`T`$ is an oriented fibre and $`\mu `$ is a geometric monodromy diffeomorphism of $`T`$ obtained by turning in the positive direction determined by the chosen orientation of $`T`$ (recall that this is precisely the point were we use the given orientation of $`M(𝒮)`$). But it is known that $`\mu `$ can be reconstructed up to isotopy by its action on $`L=H_1(T,𝐙)`$, that is, by the algebraic monodromy operator $`mSL(L)`$. Moreover, to fix an orientation of $`T`$ is the same as to fix an orientation of $`L`$. As explained in section 5.3, such an orientation can be encoded in a symplectic isomorphism $`\omega :L\stackrel{ˇ}{L}`$. Denote by $`𝒞(L,\omega ,m)`$ the cusp singularity associated to an oriented lattice $`(L,\omega )`$ and an algebraic monodromy operator $`mSL(L)`$ with $`trm3`$. If one changes the orientation of the base of the torus fibration, one gets the triple $`(L,\omega ,m^1)`$. This shows that: $$𝒞(L,\omega ,m)𝒞(L,\omega ,m^1).$$ When one changes the orientation of $`M(𝒮)`$, we see that the cone $`(L,\sigma )`$ is replaced by a supplemetary one. In view of Section 5.3, we deduce that the two cones are dual to each other. In this sense, we get the following analog of Proposition 9.3: ###### Proposition 9.6. Let $`(𝒮,0)`$ be a cusp singularity. Then the canonical plumbing structures with respect to the two possible orientations of $`M(𝒮)`$ are dual to each other. More precisely, if $`(𝒮,0)(𝒞(L,\omega ,m),0)`$, then $`M(𝒮)`$ is orientation preserving diffeomorphic to $`M(\stackrel{ˇ}{𝒮})`$, where $`(\stackrel{ˇ}{𝒮},0)(𝒞(L,\omega ,m),0)`$. ### 9.3. The other singularity boundaries As in the two previous cases, we first define the plumbing structure $`𝒫(M(𝒮))`$. Consider the canonical graph structure $`(𝒯_{can},_{can})`$ on $`M(𝒮)`$. We do our construction starting from the neighborhoods of the JSJ tori (the elements of $`𝒯_{can}`$) and the exceptional fibres in $`_{can}`$. $``$ For each component $`T`$ of $`𝒯_{can}`$, consider a saturated tubular neighborhood $`N(T)`$. We choose them pairwise disjoint. So, each manifold $`N(T)`$ is a thick torus. We consider on each one of its boundary components a fibre of $`_{can}`$. Denote these fibres by $`\gamma (T),\delta (T)`$. We consider on $`N(T)`$ the restriction of the orientation of $`M(𝒮)`$. Consider the associated Hirzebruch-Jung plumbing structure $`𝒫(N(T),\gamma (T),\delta (T))`$ (see Definition 7.17). Replace the Seifert structure on $`N(T)`$ induced from $`_{can}`$ with this plumbing structure. Then eliminate the boundary components of $`N(T)`$ from the tori present in $`M(𝒮)`$ (by construction, the fibrations coming from both sides agree on them up to isotopy). $``$ For each exceptional fibre $`F`$, consider a solid torus $`N(F)`$, which is a saturated tubular neighborhood of $`F`$. Choose those neighborhoods pairwise disjoint. On the boundary of $`N(F)`$, take a fiber $`\gamma (F)`$ of $`_{can}`$. Consider the associated Hirzebruch-Jung plumbing structure $`𝒫(N(F),\gamma (F))`$ (see Definition 7.18). Replace the Seifert structure on $`N(F)`$ induced from $`_{can}`$ with this plumbing structure. Then eliminate the boundary component of $`N(F)`$ from the tori present inside $`M(𝒮)`$ (by construction, the fibrations coming from both sides agree on it up to isotopy). Denote by $`𝒫(M(𝒮))`$ the plumbing structure constructed like this on $`M(𝒮)`$. Proof of Theorem 9.1: The proof is very similar to the ones explained in the two previous cases, but starting this time from the canonical graph structure on $`M(𝒮)`$. The main point is Proposition 8.15. We leave the details to the reader. $`\mathrm{}`$ ### 9.4. The invariance theorem Let $`(𝒮,0)`$ be a normal surface singularity. In , Neumann proved that the weighted dual graph $`\mathrm{\Gamma }(p_{mnc})`$ of the exceptional divisor of its minimal normal crossings resolution $`p_{mnc}`$ is determined by the oriented manifold $`M(𝒮)`$. But he says nothing about the action of the group $`\mathrm{Diff}^+(M(𝒮))`$ on $`(𝒯(p_{mnc}),(p_{mnc}))`$. As a corollary of Theorem 9.1 we get: ###### Theorem 9.7. The plumbing structure $`(𝒯(p_{mnc}),(p_{mnc}))`$ is invariant up to isotopy by the group $`\mathrm{Diff}^+(M(𝒮))`$. Proof: Suppose first that $`M(𝒮)`$ is not a lens space or a torus fibration. As the canonical graph structure on it is invariant by the group $`\mathrm{Diff}^+(M(𝒮))`$ up to isotopy, we deduce that the canonical plumbing structure is also invariant up to isotopy by this group. This conclusion is also true when $`M(𝒮)`$ is a lens space or a torus fibration, as one starts in the construction of $`𝒫(M(𝒮))`$ from tori which are invariant up to isotopy. Then we apply Theorem 9.1. $`\mathrm{}`$ An easy study of the fibres of $`(p_{mnc})`$ in the neighborhoods of the tori of $`𝒯(p_{mnc})`$ which correspond to self-intersection points of components of $`E_{mnc}`$ show that the analogous statement about the minimal good normal crossings resolution of $`𝒮`$ is also true. We arrived at the conclusion that the affirmation of Theorem 9.7 was true while we were thinking about the natural contact structure on $`M(𝒮)`$ (see Caubel, Némethi & Popescu-Pampu ). Indeed, in that paper we prove that for normal surface singularities, the natural contact structure depends only on the topology of $`M(𝒮)`$ up to contactomorphisms. It was then natural to look at the subgroup of $`\mathrm{Diff}^+(M(𝒮))`$ which leaves it invariant up to isotopy. Presently, we do not know how to characterize it. But we realized that the homotopy type of the underlying unoriented plane field was invariant by the full group $`\mathrm{Diff}^+(M(𝒮))`$, provided that Theorem 9.7 was true (see \[11, section 5\]).
warning/0506/cond-mat0506186.html	ar5iv	text	# Self-diffusion of rod-like viruses in the nematic phase ## 1 Introduction Suspensions of semi-flexible polymers exhibit a variety of dynamical phenomena, of great importance to both physics and biology, that are still only partially understood. Advances over the past decade include direct visual evidence for a reptation-like diffusion of individual polymers in a highly entangled isotropic solution and shape anisotropy of a single polymer . If the concentration of the polymers is increased, a suspension undergoes a first order phase transition to a nematic phase, which has long range orientational order but no long range positional order. As a result of the broken orientational symmetry it is expected that the diffusion of polymers in the nematic liquid crystals will be drastically different from that in concentrated isotropic solutions. While the static phase behavior of semi-flexible nematic polymers is well understood in terms of the Onsager theory and its extensions by Khoklov and Semenov , the dynamics of semi-flexible polymers in the nematic phase is much less explored . In this paper, we determine the concentration dependence of the anisotropic diffusion of semi-flexible viruses in a nematic phase and compare it to the diffusion in the isotropic phase. Experimentally, the only data on the translational diffusion of colloidal rods in the nematic phase was taken in a mixture of labelled and unlabelled polydysperse boehmite rods using fluorescence recovery after photobleaching (FRAP) . Theoretically, molecular dynamics simulations were performed on hard spherocylinder and ellipsoidal systems from which the anisotropic diffusion data was extracted . The anisotropic diffusion has also been studied in low molecular weight thermotropic liquid crystals using NMR spectroscopy or inelastic scattering of neutrons . Real space microscopy is a powerful method that can reveal dynamics of colloidal and polymeric liquid systems that are inaccessible to other traditional techniques . We use digital microscopy to directly visualize the dynamics of fluorescently labelled fd in a nematic background of the unlabelled fd. The advantage of this method is an easy interpretation of data and no need to obtain macroscopically aligned monodomains in magnetic fields. The advantages of using fd are its large contour length which can be easily visualized with optical microscope and its phase behavior which can be quantitatively described with the Onsager theory extended to account for electrostatic repulsive interactions and semi-flexibility . Viruses such as fd and TMV have been used earlier to study the rod dynamics in the isotropic phase . ## 2 Experiment methods The physical characteristics of the bacteriophage fd are its length L=880 nm, diameter D=6.6 nm, persistence length of 2200 nm and a surface charge of 10 e<sup>-</sup>/nm at pH 8.2 . Bacteriophage fd suspension forms isotropic, cholesteric and smectic phases with increasing concentration . The free energy difference between the cholesteric and nematic phase is very small and locally the cholesteric phase is identical to nematic. We expect that at short time scales the diffusion of the rods for these two cases would be the same. Hereafter, we refer to the liquid crystalline phase at intermediate concentration as a nematic instead of a cholesteric. The fd virus was prepared according to a standard biological protocol using XL1-Blue strain of E. coli as the host bacteria . The yields are approximately 50 mg of fd per liter of infected bacteria and virus is typically grown in 6 liter batches. Subsequently, the virus is purified by repetitive centrifugation (108,000 g for 5 hours) and re-dispersed in a 20 mM phosphate buffer at pH=7.5. First order isotropic-nematic (I-N) phase transition for fd under these conditions takes place at a rod concentration of 15.5 mg/ml. Fluorescently labelled fd viruses were prepared by mixing 1 mg of fd with 1 mg of succinimidyl ester Alexa-488 (Molecular Probes) for 1 hour. The dye reacts with free amine groups on the virus surface to form irreversible covalent bonds. The reaction is carried out in small volume (100 $`\mu `$l, 100 mM phosphate buffer, pH=8.0) to ensure a high degree of labelling. Excess dye was removed by repeated centrifugation steps. Absorbance spectroscopy indicates that there are approximately 300 dye molecules per each fd virus. Viruses labelled with fluorescein isothiocynante, a dye very similar to Alexa 488, exhibit the phase behavior identical to that of unlabelled virus. Since liquid crystalline phase behavior is a sensitive test of interaction potential, it is reasonable to assume that the interaction potential between labelled viruses is very similar to that between unlabeled viruses. The samples were prepared by mixing one unit of anti-oxygen solution (2 mg/ml glucose oxidase, 0.35 mg/ml catalase, 30 mg/ml glucose and 5% $`\beta `$-mercaptoethanol), one unit of a dilute dispersion of Alexa 488 labelled viruses and eight units of the concentrated fd virus suspension at the desired concentration. Under these conditions the fluorescently labelled viruses are relatively photostable and it is possible to continuously observe rods for 3-5 minutes without significant photobleaching. The ratio of labelled to unlabelled particles is roughly kept at 1:30000. The samples were prepared by placing 4 $`\mu `$l of solution between a No 1.5 cover slip and coverslide. The thickness of the samples is about 10 $`\mu `$m. Thin samples are important to reduce the signal of out-of-focus particles. Samples are equilibrated for half an hour, allowing flows to subside and liquid crystalline defects to anneal. We have analyzed data at various distances from the wall and have not been able to observe a significant influence of wall on the diffusion of viruses. For imaging we used an inverted Nikon TE-2000 microscope equipped with $`100\times `$ 1.4 NA PlanApo oil immersion objective, a 100 W mercury lamp and a fluorescence cube for Alexa 488 fluorescent dye. The images where taken with a cooled CCD Camera (CoolSnap HQ, Roper Scientific) set to an exposure time of 60 ms, running in a overlap mode at a rate of 16 frames per second with $`2\times 2`$ binning. The pixel size was 129 nm and the field of view was 89 $`\mu `$m $`\times `$ 66 $`\mu `$m. Typically there were around hundred fluorescently labeled rods in the field of view. For each fd concentration ten sequences of four hundred images were recorded. ## 3 Analysis method Figure 1a shows a typical image of fluorescently labelled rods in a background nematic of unlabelled rods. Due to limited spatial and temporal resolution of the optical microscope, labelled fd appear as a slightly anisotropic rod, although the actual aspect ratio is larger then 100. To measure the anisotropic diffusion in the nematic phase, it is first necessary to determine the nematic director which has to be uniform within a field of view. Spatial distortion of the nematic would significantly affects our results. The centers of mass and orientation of rods are obtained sequentially. In a first step, a smoothed image is used to identify the rods and obtain the coordinates of its center of mass using image processing code written in IDL . Subsequently, a two dimensional gaussian fit around a center of mass of each rod is performed (Fig. 1b). From this fit the orientation of each rod-like virus is obtained. This procedure is then repeated for a sequence of images. The length of a trajectory is usually limited to a few seconds, after which the particles diffuse out of focus. In Fig. 2a and b we plot the trajectories of an ensamble of particles for both isotropic and nematic sample. As expected the trajectories in the isotropic phase are spherically symmetric (Fig. 2a) while those in the nematic phase exhibit a pronounced anisotropy (Fig. 2a). The symmetric nature of the distribution indicates that there is no drift or flow in our samples. We obtain the orientation of the nematic director using two independent methods. One method is to measure the main axis of the distribution shown in Fig. 2b. This procedure assumes that the diffusion is largest along the nematic director. An alternative method is to plot a histogram of rod orientations which are obtained from 2D gaussian fits to each rod (Fig. 1b). The resulting orientational distribution function (ODF) is shown in Fig. 2c. In principle, it should be possible to obtain both the nematic director and order parameter from ODF shown Fig. 2c. We find that the order parameter obtained in such a way is systematically higher then the order parameter obtained from more reliable x-ray experiment . This is due to significant rotational diffusion each rod undergoes during an exposure time of 60 ms. The differences in the orientation of the nematic director obtained using these two methods is always less then 5 degrees. For the example shown in Figs. 2b and c, we obtain a nematic director at an angle of $`31.2^{}`$ while the peak of the orientational distribution function lies at $`30.2^{}`$. The director can be “placed” along one of the two main axis by rotating the lab-frame. The diffusion coefficients of the rods parallel ($`D_{}`$) and perpendicular ($`D_{}`$) to the director are calculated from the x’- and y’-component of the mean square displacement. When director lies along the y’-axis, $`D_{}`$ and $`D_{}`$ are given by: $`D_{}={\displaystyle \frac{1}{N}}{\displaystyle \frac{1}{2}}{\displaystyle \{y_i^{}(t)y_i^{}(0)\}^2}`$ (1) $`D_{}={\displaystyle \frac{1}{N}}{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \{x_i^{}(t)x_i^{}(0)\}^2},`$ (2) where $`N`$ is the number of traced particles. To obtain $`D_{}`$, $`D_x`$ is multiplied with $`\sqrt{2}`$ since only one component of the diffusion perpendicular to the director is measured. The underlying assumption of our analysis is that the nematic director is oriented in the field of view. For 10 $`\mu `$m thin samples this is reasonable. ## 4 Results and discussion Typical mean square displacements (MSD) are shown in Fig. 3 for samples in an isotropic and nematic phase. On average the mean square displacement was linear over fifty frames in the nematic phase, but only about twenty five frames in the isotropic phase. The diffusion perpendicular to the director is slower in the nematic phase as compared to the isotropic phase. Therefore in the nematic phase, the particles stay longer in focus and can be tracked for a longer time. Since the MSD is linear over the entire time range and displacements are up to a few times the particle length, we are measuring pure long-time self-diffusion. Visual inspection of the trajectories in the concentrated isotropic phase, just below I-N coexistence shows no characteristics of the reptation observed in suspensions of long DNA fragments or actin filaments . This points to the fact that fd is very weakly entangled in a concentrated isotropic suspension. This is in agreement with recent microrheology measurements of fd suspensions . We note that MSD’s obtained from few hundred trajectories within a single field of view are very accurate. However, if we move to another region of the sample sample we obtain MSD with slightly different slope. This leads to conclusion that the largest source of error in measuring the anisotropic diffusion coefficient is the uniformity of the nematic director within the field of view. The concentration dependence of the anisotropic diffusion constants is shown in Fig. 4a. The nematic phase melts into a isotropic phase at low concentrations and freezez into a smectic phase at high concentrations. We made an attempt to measure the diffusion of rods in the smectic phase, but have not seen any appreciable diffusion on optical length scales over a time period of minutes. The most strinking feature of our data is a strong discontinuity in the behavior of the diffusion constant at the I-N phase transition. Compared to diffusions in isotropic case $`D_{\text{iso}}`$, $`D_{}`$ is larger by a factor of four, while $`D_{}`$ is smaller by a factor of two. The concentration dependence of $`D_{}`$ and $`D_{}`$ exhibit different behavior. With increasing concentration, for $`D_{}`$ we measure an initial plateau, which is followed by a broad region where the diffusion rate decreases monotonically. $`D_{}`$, however, shows a monotonic decrease of the diffusion constant over the whole concentration range where nematic phase is stable. It is useful to compare our results to previous theoretical and experimental work, especially the measurements of the diffusion coefficient for silica coated boehmite rods . In this work authors measure $`D_{}/D_{}2`$ for monodomain nematic samples which are in coexistence with isotropic phase. This is significantly different from $`D_{}/D_{}7.5`$ for fd virus. Another significant difference is that results on boehmite indicate that both $`D_{}`$ and $`D_{}`$ are smaller then $`D_{\text{iso}}`$. In contrast to our measurements where $`D_{}`$ is much larger and $`D_{}`$ is much smaller then $`D_{\text{iso}}`$. When comparing our data to simulations of the diffusion of hard spherocylinders and ellipsoids , one needs to compare equivalent samples. Scaling to rod concentration where the I-N transition takes place would be erroneous, since fd virus is a semi-flexible rod. The semi-flexibility of the virus drives the isotropic-nematic phase transition to higher concentrations and it significantly decreases the order parameter of the nematic phase in coexistence with the isotropic phase . We choose to compare data and simulations at the same value of the nematic order parameter which is determined independently . For fd, the nematic order parameter is 0.65 at the I-N coexistence, monotonically increases with increasing rod concentration and saturates at high rod concentration. Experiment and simulation qualitatively agree and both show a rapid increase of $`D_{}/D_{}`$ ratio with increasing nematic order parameter (Fig. 4b). We note that there is a discrepancy between the simulations results obtain in references which might be due to different systems studied in these two paper. Interestingly, simulations predict that upon increasing rod concentration beyond I-N coexistence $`D_{}`$ increases and subsequently upon approaching the smectic phase it will decrease. The authors argues that the non-monotonic behavior of $`D_{}`$ is the result of the interplay between two effects. First, with increasing rod concentration the nematic order parameter increases which enhances $`D_{}`$. Second, with increasing rod concentration there is less free volume which leads to decrease of $`D_{}`$. The author further argues that the first effect dominates at low rod concentrations where the nematic order parameter rapidly increases while the second effect dominates at high rod concentrations where the nematic order parameter is almost saturated. In contrast, both of these effects contribute to a monotonic decrease in $`D_{}`$ with increasing concentration, which is observed in simulations. Due to relatively large error in our experimental data, it is not clear if the behavior of $`D_{}`$ is non-monotonic. There is an initial hesitation, but $`D_{}`$ decreases over most of the concentration range. This difference between simulations and experiment might be because we compare experiments of semi-flexible fd to simulations of perfectly rigid rods. Compared to semi-flexible rods, the order parameter of rigid rods increases much faster with increasing rod concentration . It would be of interest to extend our measurements to rotational diffusion in the isotropic and nematic phase. At present the rod undergoes significant rotational diffusion during each exposure which reduces resolution and prevents accurate determination of the instantaneous orientation of a rod. It might be possible to significantly reduce the exposure time by either using a more sensitive CCD camera or a more intense laser as a illumination source. ## 5 Conclusions Using fluorescence microscopy we have visualized rod-like viruses and measured the anisotropic long-time self-diffusion coefficients in the isotropic and nematic phase. In the nematic phase the diffusion along the director and the diffusion perpendicular to the director decreases monotonically with increasing rod concentration. The ratio of parallel to perpendicular diffusion increases monotonically with increasing rod concentration. The results compare qualitatively with simulations on hard rods with moderate aspect ratios. ###### Acknowledgements. Pavlik Lettinga is supported in part by Transregio SFB TR6, ”Physics of colloidal dispersions in external fields”. Zvonimir Dogic is supported by Junior Fellowship from Rowland Institute at Harvard.
warning/0506/physics0506187.html	ar5iv	text	# Attosecond electron thermalization by laser-driven electron recollision in atoms ## Abstract Nonsequential multiple ionization of atoms in intense laser fields is initiated by a recollision between an electron, freed by tunneling, and its parent ion. Following recollision, the initial electron shares its energy with several bound electrons. We use a classical model based on rapid electron thermalization to interpret recent experiments. For neon, good agreement with the available data is obtained with an upper bound of 460 attoseconds for the thermalization time. Atoms exposed to intense laser fields ionize. The freed electron and its ionic partner are accelerated by the laser field away from each other. When the field changes sign, it may drive the electron into a recollision with the ion. This simple mechanism, which for high-intensity low-frequency fields is largely classical, governs many laser-atom processes such as high-order harmonic generation (HHG), high-order above-threshold ionization (HATI), and nonsequential double ionization (NSDI) and explains the gross features of the spectra observed corkum . The period of the commonly applied titanium-sapphire (Ti:Sa) laser is about 2.7 fs. The recollision physics unfold on the time scale of a small fraction of the laser period. Therefore, the analysis of laser-induced recollision phenomena provides access to the inner-atomic dynamics on the attosecond time scale and, indeed, the focus of recent investigations has moved into this temporal domain. For example, it has brought molecular imaging with subangstrom spatial and subfemtosecond temporal resolution within reach Niikura ; molimag . The advent of phase-stabilized infrared few-cycle pulses and of uv pulses of attosecond duration allows even more control Baltuska , but neither are necessary for a study of the attosecond dynamics. In this Letter, we analyze recent experiments on nonsequential multiple ionization (NSMI) of neon in which the momentum distributions of Ne<sup>N+</sup> $`(N=3,4)`$ were measured at two intensities. We use the fact that the time-dependent laser field provides a clock — the field accelerates the ion to a final velocity that depends on the times at which the $`N`$ electrons ionized. This “streak camera” streak therefore measures the range of times of ionization. Comparing the measured momentum distributions to those predicted by a classical model, we infer that the recolliding electron thermalizes with the $`N1`$ bound electrons in less than 500 attoseconds. To our knowledge, no other method allows such a low upper bound for thermalization times within atoms to be measured. Nonsequential double and multiple ionization is defined by the fact that it is not sequential, that is, it is not the product of a sequence of uncorrelated single-ionization events. NSDI and NSMI require electron-electron correlation as a necessary precondition Ffm2000 ; FrMBI2000 . Even for the very simplest such process – NSDI of helium – a fully quantum-mechanical description from first principles, i.e. by solution of the time-dependent Schrödinger equation in six spatial dimensions, has not been accomplished yet taylor , and for the heavier atoms it is clearly out of the question. This leaves approximate quantum-mechanical approaches, such as density-functional methods bauer , $`S`$-matrix methods that try to identify the most relevant terms of an appropriate perturbative expansion faisal ; KBRS , or classical-trajectory methods china ; eberly . For multiple (triple and higher) ionization, which occurs under the same conditions as NSDI if the laser intensity is high enough, any description from first principles appears to be utterly out of reach. It seems equally hopeless, for three electrons and more, to identify the relevant diagrams in the microscopic $`S`$-matrix approach. The model we here propose is in the spirit of the one of Ref. KBRS , but essentially classical. We investigate the scenario wherein NSMI is effected by one single recollision. We assume that the pertinent electron tunnels into the continuum with zero velocity at the ionization time $`t^{}`$ according to the time-dependent rate $`R(t^{})`$, for which we adopt the standard quasi-static rate LL . Thereafter, we turn to an entirely classical description: The laser field may drive the electron back to its parent ion at a later time $`t`$, which is a function $`t(t^{})`$ of the ionization time and can be easily evaluated. We assume that the energy $`E_{\mathrm{ret}}(t)`$ of the returning electron be completely thermalized among the ensemble of participating electrons, that is, the returning electron and the $`N1`$ electrons to be freed. These $`N`$ electrons then form an excited complex with the total energy (with respect to the continuum threshold) $`E_{\mathrm{ret}}(t)E_0^{(N)}`$, where $`E_0^{(N)}>0`$ denotes the total ionization potential of the $`N1`$ (up to the recollision time $`t`$ inactive) electrons. The distribution of energy and momentum over the $`N`$ electrons is assumed to be completely statistical and only governed by the available phase space. At the time $`t+\mathrm{\Delta }t`$, the $`N`$ electrons become free to move in the laser field, which is described by the vector potential $`𝐀(t)`$ such that $`𝐀(t)=\mathrm{𝟎}`$ outside the pulse. The corresponding distribution of the final electron momenta $`𝐩_n(n=1,\mathrm{},N)`$ is proportional to $`F(𝐩_1,𝐩_2,\mathrm{},𝐩_N)={\displaystyle }dt^{}R(t^{})\delta (E_0^{(N)}E_{\mathrm{ret}}(t)`$ $`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}[𝐩_n+𝐀(t+\mathrm{\Delta }t)]^2),`$ (1) where the integral extends over the ionization time $`t^{}`$. The $`\delta `$ function expresses the fact that the total kinetic energy of the $`N`$ participating electrons at the time $`t+\mathrm{\Delta }t`$ is fixed by the first-ionized electron at its recollision time $`t`$. This constitutes the one and only condition on the final momenta $`𝐩_n`$. The only free parameter of this model is the time delay $`\mathrm{\Delta }t`$ between the recollision time and the time when the electrons become free. It is the sum of a thermalization time $`\mathrm{\Delta }t_{\mathrm{th}}`$ – the time it takes to establish the statistical ensemble – and a possible additional “dwell time”, until the electrons become free. By comparing the predictions of the model with the data, we will be able to infer a value of $`\mathrm{\Delta }t`$, which in turn provides an upper limit for the thermalization time $`\mathrm{\Delta }t_{\mathrm{th}}`$. This model is an extension to NSMI of a classical model introduced for NSDI for $`\mathrm{\Delta }t=0`$ FFetal04R ; FFetal04 . Sufficiently high above threshold, it produced momentum distributions that were virtually indistinguishable from their quantum-mechanical counterparts. Statistical models similar to the one above have been used in many areas of physics. For example, the statistical Rice-Ramsperger-Kassel-Marcus (RRKM) theory RRKM describes thermalization of molecular vibrational degrees of freedom, and for high-energy collisions of elementary particles and heavy ions statistical models have been utilized to predict the momentum spectra of the reaction products hagedorn . An excited complex as the doorway to NSMI was also considered in Ref. sachaeckh . A convenient feature of the ansatz (1) is that integration over unobserved momentum components is easily carried out. To this end, we exponentialize the $`\delta `$ function in Eq. (1) with the help of its Fourier representation $$\delta (x)=_{\mathrm{}}^{\mathrm{}}\frac{d\lambda }{2\pi }\mathrm{exp}(i\lambda x).$$ (2) Infinite integrations over the momenta $`𝐩_n`$ can then be done by Gaussian quadrature. The remaining integration over the variable $`\lambda `$ is taken care of by the formula GR $$_{\mathrm{}}^{\mathrm{}}\frac{d\lambda }{(i\lambda +ϵ)^\nu }e^{ip\lambda }=\frac{2\pi }{\mathrm{\Gamma }(\nu )}p_+^{\nu 1},$$ (3) where $`x_+^\nu =x^\nu \theta (x)`$, with $`\theta (x)`$ the unit step function and $`ϵ+0`$. For comparison with the experiments FrMBI2000 ; multi , we calculate the distribution of the momentum $`𝐏`$ of the ion. Provided the momentum of the absorbed laser photons can be neglected, momentum conservation implies $`𝐏=_{n=1}^N𝐩_n`$, so that $`F_{\mathrm{ion}}(𝐏){\displaystyle \underset{n=1}{\overset{N}{}}d^3𝐩_n\delta \left(𝐏+\underset{n=1}{\overset{N}{}}𝐩_n\right)F(𝐩_1,𝐩_2,\mathrm{},𝐩_N)}`$ $`={\displaystyle \frac{(2\pi )^{\frac{3N}{2}\frac{3}{2}}}{N^{3/2}\mathrm{\Gamma }(\frac{3}{2}(N1))}}{\displaystyle 𝑑t^{}R(t^{})\left(\mathrm{\Delta }E_{N,\mathrm{ion}}\right)_+^{\frac{3N}{2}\frac{5}{2}}}`$ (4) with $`\mathrm{\Delta }E_{N,\mathrm{ion}}E_{\mathrm{ret}}(t)E_0^{(N)}\frac{1}{2N}[𝐏N𝐀(t+\mathrm{\Delta }t)]^2`$. In the experiments thus far, the ion momentum transverse to the direction of the laser polarization is entirely or partly integrated over. In the first case, the remaining distribution of the longitudinal ion momentum $`P_{}`$ is $`F_{\mathrm{ion}}(P_{}){\displaystyle d^2𝐏_{}F_{\mathrm{ion}}(𝐏)}`$ $`={\displaystyle \frac{(2\pi )^{\frac{3N}{2}\frac{1}{2}}}{\sqrt{N}\mathrm{\Gamma }((3N1)/2)}}{\displaystyle 𝑑t^{}R(t^{})\left(\mathrm{\Delta }E_{\mathrm{ion}}\right)_+^{\frac{3N}{2}\frac{3}{2}}},`$ (5) where now $`\mathrm{\Delta }E_{\mathrm{ion}}E_{\mathrm{ret}}(t)E_0^{(N)}\frac{1}{2N}[P_{}NA(t+\mathrm{\Delta }t)]^2`$. If just one transverse-momentum component ($`P_{,2}`$, say) is integrated while the other one ($`P_{,1}P_{}`$) is observed, the corresponding distribution is $`F_{\mathrm{ion}}(P_{},P_{}){\displaystyle 𝑑P_{,2}F_{\mathrm{ion}}(𝐏)}`$ $`={\displaystyle \frac{(2\pi )^{\frac{3N}{2}1}}{N\mathrm{\Gamma }(3N/21)}}{\displaystyle 𝑑t^{}R(t^{})\left(\mathrm{\Delta }E_{N,\mathrm{ion}}\right)_+^{\frac{3N}{2}2}}`$ (6) with $`\mathrm{\Delta }E_{N,\mathrm{ion}}\mathrm{\Delta }E_{\mathrm{ion}}\frac{1}{2N}P_{}^2`$. In Fig. 1 we present calculations of the ion-momentum distributions for triple (upper panel) and quadruple (lower panel) NSMI of neon according to Eq. (5), for various values of $`\mathrm{\Delta }t`$ between 0 and $`0.2T`$. The parameters are for the experimental data of neon presented in Fig. 2 of Ref. multi . The figure displays the double-hump structure of the ion momentum in NSMI of neon. When the time delay $`\mathrm{\Delta }t`$ increases from 0, initially, the effect on the momentum distribution is small. Later, however, the center positions of the humps start moving towards zero momentum and the widths of the humps increase until the two humps begin to merge. This behavior can easily be understood from the recollision kinematics, which are illustrated in Fig. 2. The final electrons undergo maximal acceleration if they are released near a zero crossing of the electric field, which in the figure occurs at $`t=T`$. For the example of triple ionization of neon, the earliest recollision with $`E_{\mathrm{ret}}>E_0^{(3)}`$ takes place at $`t=0.74T`$. Already with a delay of $`\mathrm{\Delta }t>0.26T`$, all electrons will be released after the zero crossing, which results in significantly lower ion momenta. In Fig. 3 we compare the results of the statistical model with the data of Ref. multi . The intensities are those given for the experiment. We display momentum distributions calculated from Eq. (5) for zero delay and for $`\mathrm{\Delta }t=0.17T`$. The latter value was chosen to yield optimal agreement for the entire set of data. We notice, in particular, that with this nonzero delay the model reproduces the maxima of the experimental ion-momentum distribution. This removes a longstanding discrepancy between models of the type discussed in Refs. FFetal04R ; FFetal04 ; KBRS and the data. For triple ionization, the calculated momentum distributions are wider than those of the data, in particular for the higher intensity, even though the data may include a contribution of the partly sequential channel Ne $``$ Ne$`{}_{}{}^{+}`$ Ne<sup>3+</sup>, or the contribution of a recollision-excitation channel, both of which are, of course, not part of the model. Figure 4 exhibits the results of our thermalization model for $`N=3`$ for the conditions of Ref. FrMBI2000 . The distribution of two components of the ion momentum is presented, the component parallel to the laser field and one transverse component, which therefore provides a more stringent test of the model. The third component is integrated over in the data, which corresponds to Eq. (6). The intensity given in the experiment is 1.5 PWcm<sup>-2</sup>; we obtain good agreement with the data for the lower intensity of 1.0 PWcm<sup>-2</sup> and $`\mathrm{\Delta }t=0.17T`$ as before footnote . The nonzero delay $`\mathrm{\Delta }t`$ has little effect on the transverse width of the distribution, but it causes an elongation in the longitudinal direction and moves the centers to lower momenta, markedly improving the agreement with the data. Encouraged by the good agreement between the model and the data, we interpret the delay for which we observed optimal agreement as an upper bound of the thermalization time $`\mathrm{\Delta }t_{\mathrm{th}}`$, as argued above. This yields $`\mathrm{\Delta }t_{\mathrm{th}}<\mathrm{\Delta }t_{\mathrm{opt}}=0.17T460`$ as. A lower limit of the thermalization time should be given by the inverse of the plasma frequency for an electron density $`\rho _\mathrm{e}=N`$ in atomic units. This produces $`\mathrm{\Delta }t_{\mathrm{th}}<\sqrt{\pi /N}`$, which is of the order of the atomic unit of time. All of the above discussion has been for neon. NSDI of argon appears to be governed by a recollision-excitation scenario ArvsNe . We note in passing that for (quadruple) NSMI of argon we get good agreement of our statistical thermalization model with the data multi for the experimental intensity and $`\mathrm{\Delta }t=0.35T`$, twice as long as for neon. Details will be given elsewhere. To test the model further, and possibly to set a tighter upper limit on the thermalization time, it is necessary to restrict the time range $`\mathrm{\Delta }t_{\mathrm{CM}}`$ of recollision (Fig. 2). If thermalization is more rapid than 460 as, then the width of the ion-momentum distribution will decrease as we restrict $`\mathrm{\Delta }t_{\mathrm{rec}}`$. There are a number of ways to minimize $`\mathrm{\Delta }t_{\mathrm{rec}}`$. Within limits, all that is needed is to lower the light intensity or to increase the laser frequency. However, the best way to minimize $`\mathrm{\Delta }t_{\mathrm{rec}}`$ and to control the time of recollision is to use a second harmonic field, polarized perpendicular to the fundamental. The combined requirement that the electron and ion recollide in both directions allows the time of recollision to be precisely determined and controlled via the relative phase of the two beams. The limit on the thermalization time that we determine (as well as the much tighter bounds that seem feasible in future experiments) should apply to stationary electron-atom scattering in general. From a collision-physics perspective, as a result of streaking, the time-dependent laser field reveals information that is hard to obtain by other means. The streaking principle should also be applicable for studying nuclear dynamics. As in NSMI, nuclear processes also can be initiated by laser-controlled re-collision nuclear (of course, at much higher intensity). Any nuclear decay process that results in a mass or charge change of the fragments will be streaked by the laser field just as electrons and ions are streaked in our case. To summarize, by comparison of experimental data for triple and quadruple nonsequential ionization of neon with a simple statistical recollision model where the returning electron thermalizes with a subset of the bound electrons, we have been able to conclude that (i) for neon such a model appears to contain the most relevant physics, and (ii) the time for this thermalization to occur is extremely fast, well below one femtosecond. We gratefully acknowledge stimulating discussions with G.G. Paulus, H. Rottke, and W. Sandner.
warning/0506/hep-ph0506315.html	ar5iv	text	# Next-to-leading order QCD corrections to single-inclusive hadron production in transversely polarized 𝐩𝐩 and 𝐩̄⁢𝐩 collisions ## I Introduction The partonic structure of spin-1/2 targets at the leading-twist level is characterized entirely by the unpolarized, longitudinally polarized, and transversely polarized distribution functions $`f`$, $`\mathrm{\Delta }f`$, and $`\delta f`$, respectively ref:jaffeji . Of these, the “transversity” distributions $`\delta f`$ remain virtually unknown. They are defined ref:jaffeji ; ref:ralston ; ref:artru ; ref:ratcliffe as the differences of probabilities for finding a parton of flavor $`f`$ at scale $`\mu `$ and light-cone momentum fraction $`x`$ with its spin aligned ($``$) or anti-aligned ($``$) with that of the transversely polarized nucleon: $$\delta f(x,\mu )f_{}(x,\mu )f_{}(x,\mu ).$$ (1) A program of polarized $`pp`$ collisions is now underway at the BNL Relativistic Heavy Ion Collider (RHIC) ref:rhic , aiming at further unraveling the spin structure of the proton. Collisions of transversely polarized protons are hoped to give information on transversity through, e.g., the measurement of double-spin asymmetries $$A_{\mathrm{TT}}=\frac{\frac{1}{2}[d\sigma ()d\sigma ()]}{\frac{1}{2}[d\sigma ()+d\sigma ()]}\frac{d\delta \sigma }{d\sigma }$$ (2) for various reactions with observed produced high-transverse momentum ($`p_T`$) or invariant mass. The best-studied, and perhaps most promising among these, is the Drell-Yan process ref:ralston ; ref:drellyan ; ref:drellyan2 , which offers the largest spin asymmetries but whose main drawback is the rather moderate event rate. Other reactions, such as high-$`p_T`$ prompt-photon, pion, or jet production, are much more copious, but suffer from fairly small spin asymmetries ref:artru ; ref:jaffesaito ; ref:ji92 ; ref:attlo ; ref:photonnlo , due to large contributions from gluon-gluon and quark-gluon scattering only present in the unpolarized cross section in the denominator of Eq. (2). Very recently, it has also been proposed to extract transversity from measurements of $`A_{\mathrm{TT}}`$ in transversely polarized $`\overline{p}p`$ collisions at the planned GSI-FAIR facility ref:pax ; ref:assia ; ref:anselmino near Darmstadt, Germany. For later stages of operations, there are plans to have an asymmetric $`\overline{p}p`$ collider, with moderate proton and antiproton energies of 3.5 and 15 GeV, respectively. So far, theoretical work has focused on the Drell-Yan process. It was found that the expected spin asymmetries could be very large, possibly reaching several tens of per cents ref:anselmino ; ref:radici ; ref:resum . This can be readily understood, because for the GSI kinematics only partons with rather large momentum fractions scatter off each other, and in $`\overline{p}p`$ collisions the relevant lowest order (LO) process, $`q\overline{q}`$ annihilation, will receive large contributions from valence quarks, which are expected to carry strong polarization. This appears to make the proposed measurements at GSI particularly interesting for learning about transversity. The theoretical framework for GSI kinematics is somewhat more involved than for RHIC, since perturbative all-order resummations of large logarithmic contributions to the partonic cross sections are particularly important. For the Drell-Yan process, these have been addressed in detail recently in ref:resum . In any case, information from RHIC and from GSI experiments will be complementary, due to the very different kinematics accessed, and very likely both will be needed to gain sufficient knowledge about the $`\delta f`$ over a large range in $`x`$. In this paper, we perform a detailed study of high-$`p_T`$ single-inclusive pion production in transversely polarized $`pp`$ and $`\overline{p}p`$ collisions. In particular, we derive the next-to-leading order (NLO) QCD corrections to the relevant partonic cross sections. In general, these are indispensable for arriving at a firmer theoretical framework for analyzing experimental data in terms of parton densities. For the calculation we employ a recently developed “projection technique” for treating the phase space integrals in the presence of the $`\mathrm{cos}(2\mathrm{\Phi })`$ azimuthal-angular dependence associated with transverse polarization ref:photonnlo . We will apply our analytical results in phenomenological studies for $`pp`$ and $`\overline{p}p`$ scattering at RHIC and GSI energies, respectively. Regarding the latter, to our knowledge our study is the first to propose accessing transversity via the process $`\overline{p}p\pi X`$. We hope that such measurements would be possible with the proposed PAX ref:pax and ASSIA ref:assia experiments, but this reaction should also be of great interest for the PANDA Collaboration ref:panda . As we shall see, the spin asymmetry for $`\overline{p}p\pi X`$ is expected to be much smaller than that for the Drell-Yan process, a drawback that may be compensated for by the much higher event rates and, therefore, the much better statistical accuracy. We will also find, however, that the size of the NLO corrections and the scale dependence are very significant at GSI energies, so that further theoretical work will be needed before one can be confident that high-$`p_T`$ pion production may be a useful tool to learn about transversity. If so, combined information from Drell-Yan and from $`\overline{p}p\pi X`$ (or other produced hadrons) could be useful since the processes probe different combinations of the transversity densities. In the next section we will very briefly review the necessary technical framework for the computation of the NLO QCD corrections. Details on the projection technique, which is a crucial tool for the calculation, can be found in Ref. ref:photonnlo where it was applied to prompt-photon production. In Sec. III we will present phenomenological results for RHIC and GSI energies. We conclude in Sec. IV. ## II Technical Framework According to the factorization theorem ref:fact , the fully differential, transverse-spin dependent, single-inclusive cross section for the reaction $`AB\pi X`$ for the production of a pion (or any other hadron) with transverse momentum $`p_T`$, azimuthal angle $`\mathrm{\Phi }`$ with respect to the initial spin axis, and pseudorapidity $`\eta `$ reads at NLO accuracy $`{\displaystyle \frac{d^3\delta \sigma }{dp_Td\eta d\mathrm{\Phi }}}`$ $`=`$ $`{\displaystyle \frac{p_T}{\pi S}}{\displaystyle \underset{abcX}{}}{\displaystyle _{1V+VW}^1}{\displaystyle \frac{dz_c}{z_c^2}}{\displaystyle _{VW/z_c}^{1(1V)/z_c}}{\displaystyle \frac{dv}{v(1v)}}{\displaystyle _{VW/vz_c}^1}{\displaystyle \frac{dw}{w}}\delta f_a(x_a,\mu )\delta f_b(x_b,\mu )D_c^\pi (z_c,\mu )`$ (3) $`\times `$ $`\left[{\displaystyle \frac{d\delta \widehat{\sigma }_{abcX}^{(0)}(v)}{dvd\mathrm{\Phi }}}\delta (1w)+{\displaystyle \frac{\alpha _s(\mu )}{\pi }}{\displaystyle \frac{d\delta \widehat{\sigma }_{abcX}^{(1)}(s,v,w,\mu )}{dvdwd\mathrm{\Phi }}}\right],`$ where the sum is over all contributing partonic channels $`abcX`$, $`AB=pp`$ or $`\overline{p}p`$, and with hadron-level variables $`V`$ $``$ $`1+{\displaystyle \frac{T}{S}},W{\displaystyle \frac{U}{S+T}},S(P_A+P_B)^2,`$ $`T`$ $``$ $`(P_AP_\pi )^2,U(P_BP_\pi )^2,`$ (4) in obvious notation of the momenta. The corresponding partonic quantities are given by $`v`$ $``$ $`1+{\displaystyle \frac{t}{s}},w{\displaystyle \frac{u}{s+t}},s(p_a+p_b)^2,`$ $`t`$ $``$ $`(p_ap_c)^2,u(p_bp_c)^2.`$ (5) Neglecting all masses, one has the relations $`s`$ $`=`$ $`x_ax_bS,t={\displaystyle \frac{x_a}{z_c}}T,u={\displaystyle \frac{x_b}{z_c}}U,`$ $`x_a`$ $`=`$ $`{\displaystyle \frac{VW}{vwz_c}},x_b={\displaystyle \frac{1V}{(1v)z_c}}.`$ (6) The transversity densities in Eq. (3) always refer to those for a parent proton even for $`AB=\overline{p}p`$, i.e., we use the charge conjugation property $`\delta f_a^{\overline{p}}=\delta f_{\overline{a}}^p`$. The fact that we are observing a specific hadron in the reaction requires the introduction of additional long-distance functions in Eq. (3), the parton-to-pion fragmentation functions $`D_c^\pi `$. The $`d\delta \widehat{\sigma }_{abcX}^{(i)}`$ are the LO ($`i=0`$) and NLO ($`i=1`$) contributions in the cross sections for the partonic reactions $`abcX`$. Finally, $`\mu `$ collectively denotes the renormalization and factorization scales, which we will always take as equal for simplicity. We now give a few technical details of the NLO calculation. Projection on a definite polarization state for the initial partons involves the Dirac matrix $`\gamma _5`$. It is well known that in dimensional regularization, which we will use to regularize the ultraviolet, infrared, and collinear singularities at intermediate stages of the calculation, the treatment of $`\gamma _5`$ is in general a subtle issue. However, owing to the chirally odd nature of transversity, in our calculation all Dirac traces contain two $`\gamma _5`$ matrices, and, therefore, using the “HVBM scheme” ref:hvbm or a naive, totally anticommuting $`\gamma _5`$ in $`n4`$ dimensions must give the same results which is also a useful check for the correctness of the calculation. The transverse polarization vectors of the initial hadrons give rise to a characteristic dependence of the cross section on the azimuthal angle $`\mathrm{\Phi }`$ of the observed particle. In the hadronic center-of-mass system (c.m.s.) frame, taking the initial hadrons along the $`\pm z`$ axis and their spin vectors in $`\pm x`$ direction, the $`\mathrm{\Phi }`$-dependence is of the form $`\mathrm{cos}(2\mathrm{\Phi })`$. Integration over $`\mathrm{\Phi }`$ is therefore not appropriate. Keeping $`\mathrm{\Phi }`$ fixed in the NLO calculation is however very cumbersome since standard techniques developed in the literature for performing NLO phase-space integrations rely on the choice of particular reference frames different from the one specified above. In ref:photonnlo we developed a general projection method that involves integration over all $`\mathrm{\Phi }`$, thereby allowing to keep the benefits of the standard phase space integration techniques. The trick, and the virtue of our method, is to project out the dependence of the matrix elements on the spin vectors in a covariant way, by multiplying with a covariant expression for the $`\mathrm{cos}(2\mathrm{\Phi })`$ term, and to then carry out the complete phase space integrals. To be more specific, we note that because of the $`\mathrm{cos}(2\mathrm{\Phi })`$ dependence we have the identity $$\frac{d^3\delta \sigma }{dp_Td\eta d\mathrm{\Phi }}\mathrm{cos}(2\mathrm{\Phi })_0^{2\pi }𝑑\mathrm{\Phi }^{}\frac{\mathrm{cos}(2\mathrm{\Phi }^{})}{\pi }\frac{d^3\delta \sigma }{dp_Td\eta d\mathrm{\Phi }^{}}.$$ (7) The $`\mathrm{cos}(2\mathrm{\Phi })`$ dependence actually arises through the covariant expression $$(p_c,s_a,s_b)=\frac{s}{tu}\left[2(p_cs_a)(p_cs_b)+\frac{tu}{s}(s_as_b)\right],$$ (8) where the $`s_i`$ ($`i=a,b`$) are the initial transverse spin vectors which satisfy $`s_ip_a=s_ip_b=0`$ and $`s_a^2=s_b^2=1`$. $`(p_c,s_a,s_b)`$ reduces to $`\mathrm{cos}(2\mathrm{\Phi })`$ in the hadronic c.m.s. frame. We may, therefore, use $`(p_c,s_a,s_b)/\pi `$ instead of the explicit $`\mathrm{cos}(2\mathrm{\Phi })/\pi `$ in the “projector” in the integrand of Eq. (7). For any contributing partonic channel we multiply the squared matrix element for transversely polarized initial partons, $`\delta \|M\|_{abcX}^2`$, by $`(p_c,s_a,s_b)/\pi `$. The resulting expression may then be integrated over the full azimuthal phase space in a covariant way without producing a vanishing result, unlike the case of $`\delta \|M\|^2`$ itself; see Ref. ref:photonnlo for further details. It is crucial here that the other observed (“fixed”) quantities, the hadron’s transverse momentum $`p_T`$ and pseudorapidity $`\eta `$, are determined entirely by scalar products $`(p_ap_c)`$ and $`(p_bp_c)`$, independently of the spin vectors $`s_{a,b}`$. Our method becomes particularly convenient for treating the $`23`$ scattering contributions arising at NLO where one has an additional phase space integral over the second unobserved parton in the final-state. After applying the projection method we can perform all phase space integrations by employing techniques familiar from the corresponding calculations in the unpolarized and longitudinally polarized cases ref:nlostuff ; ref:nlopion . We note that as a non-trivial check on our calculation we have also integrated all squared matrix elements over the spin vectors without using any projector at all. This amounts to integrating $`\mathrm{cos}(2\mathrm{\Phi })`$ over all $`0\mathrm{\Phi }2\pi `$, and, as expected, the final answer is zero. The use of dimensional regularization is straightforward in all this. Ultraviolet poles in the virtual diagrams are removed by the renormalization of the strong coupling constant. Infrared singularities cancel in the sum between virtual and real-emission diagrams. After this cancellation, only collinear poles are left. From the factorization theorem it follows that these need to be factored into the parton distribution and fragmentation functions. This is a standard procedure which we have also described in quite some detail in ref:nlopion ; ref:photonnlo . We use the $`\overline{\mathrm{MS}}`$ scheme throughout. After factorization, we arrive at the final result, the finite partonic NLO hard scattering cross sections. There are all in all five subprocesses that contribute for transverse polarization: $`qq`$ $``$ $`qX,`$ $`q\overline{q}`$ $``$ $`qX,`$ $`q\overline{q}`$ $``$ $`q^{}X,`$ $`q\overline{q}`$ $``$ $`gX,`$ $`qq`$ $``$ $`gX,`$ (9) where at NLO in each case $`X`$ denotes a one- or two-parton final state, summed over all possibilities and integrated over its phase space. The first four of these reactions are present at LO already. The corresponding LO transversity cross sections may be found in ref:jaffesaito ; ref:ji92 ; ref:attlo ; ref:attold . The last subprocess appears for the first time at NLO. For each of the five subprocesses, the NLO expression for the transversely polarized cross section can be cast into the following form: $`s{\displaystyle \frac{d\delta \widehat{\sigma }_{abcX}^{(1)}(s,v,w,\mu )}{dvdwd\mathrm{\Phi }}}=\mathrm{cos}(2\mathrm{\Phi })\left({\displaystyle \frac{\alpha _s(\mu )}{\pi }}\right)^2[(A_0\delta (1w)+B_0{\displaystyle \frac{1}{(1w)_+}}+C_0)\mathrm{ln}{\displaystyle \frac{\mu ^2}{s}}`$ $`+A\delta (1w)+B{\displaystyle \frac{1}{(1w)_+}}+C+D\left({\displaystyle \frac{\mathrm{ln}(1w)}{1w}}\right)_++E\mathrm{ln}w+F\mathrm{ln}v+G\mathrm{ln}(1v)`$ $`+H\mathrm{ln}(1w)+I\mathrm{ln}(1vw)+J\mathrm{ln}(1v+vw)+K{\displaystyle \frac{\mathrm{ln}w}{1w}}+L{\displaystyle \frac{\mathrm{ln}\frac{1v}{1vw}}{1w}}+M{\displaystyle \frac{\mathrm{ln}(1v+vw)}{1w}}],`$ (10) where the “plus”-distribution is defined in the usual way over the interval $`[0,1]`$. All coefficients in Eq. (10) are functions of $`v`$ and $`w`$, except those multiplying the distributions $`\delta (1w)`$, $`1/(1w)_+`$, $`\left[\mathrm{ln}(1w)/(1w)\right]_+`$ which may be written as functions just of $`v`$. Terms with distributions are present only for the subprocesses that already contribute at the Born level. The coefficients are available upon request as a Fortran code from the authors. ## III Phenomenological Results We now present some phenomenological results for single-inclusive pion production in transversely polarized $`pp`$ collisions at RHIC ($`\sqrt{S}=200`$ and 500 GeV) and asymmetric $`\overline{p}p`$ collisions at the planned GSI-FAIR facility with proton and antiproton energies of $`3.5`$ GeV and 15 GeV, respectively. Since nothing is known experimentally about transversity so far, we need to model the $`\delta f`$ for our study. Guidance is provided by the Soffer inequality ref:soffer $$2\left\|\delta q(x)\right\|q(x)+\mathrm{\Delta }q(x)$$ (11) which gives bounds for each $`\delta f`$. As in ref:attlo ; ref:photonnlo we utilize this inequality by saturating the bound at some low input scale $`\mu _00.6\mathrm{GeV}`$, choosing all signs to be positive, and using the NLO (LO) GRV ref:grv and GRSV (“standard scenario”) ref:grsv densities $`q(x,\mu _0)`$ and $`\mathrm{\Delta }q(x,\mu _0)`$, respectively. For $`\mu >\mu _0`$ the transversity densities $`\delta f(x,\mu )`$ are then obtained by evolving them at LO or NLO. We refer the reader to ref:drellyan2 ; ref:attlo for more details on our model distributions. We note that we will always perform the NLO (LO) calculations using NLO (LO) parton distribution functions and the two-loop (one-loop) expression for $`\alpha _s`$. We use the pion fragmentation functions of Ref. ref:frag which has both a LO and an NLO set. They provide a very good description of the recent RHIC data on unpolarized neutral-pion production ref:rhicdata . Figure 1 shows our estimates for the transversely polarized single-inclusive pion production cross sections at LO and NLO for the two different c.m.s. energies at RHIC. We have integrated over the range $`\|\eta \|0.38`$ in pseudorapidity, appropriate for measurements with the PHENIX detector. Since only half of the pion’s azimuthal angle is covered, we integrate over the two quadrants $`\pi /4<\mathrm{\Phi }<\pi /4`$ and $`3\pi /4<\mathrm{\Phi }<5\pi /4`$, which gives $`\left(_{\pi /4}^{\pi /4}+_{3\pi /4}^{5\pi /4}\right)\mathrm{cos}(2\mathrm{\Phi })d\mathrm{\Phi }=2`$. We have also varied simultaneously the factorization/renormalization scales $`\mu `$ in Eq. (3) within $`p_T\mu 4p_T`$; a significant decrease of scale dependence is observed when going from LO to NLO. The lower part of the Figure 1 displays the so-called “$`K`$-factor”, defined as usual as the ratio of the NLO to the LO cross section, for the scale choice $`\mu =2p_T`$. Except for small $`p_T`$, where the NLO corrections lead to a significant reduction of the cross section, the $`K`$-factor turns out to be rather moderate and close to unity. It is known that the $`K`$-factor for the unpolarized cross section is significantly larger than one at RHIC energies, see, e.g., Fig. 4 in Ref. ref:nlopion for $`\sqrt{S}=200\mathrm{GeV}`$, mostly because of large corrections found for gluon-initiated partonic channels. Therefore, one expects that the double-spin asymmetry $`A_{\mathrm{TT}}`$ at RHIC will decrease when going from LO to NLO. Indeed, as Fig. 2 shows, this is the case. Here we used the CTEQ6M (CTEQ6L1) ref:cteq6 set of unpolarized parton distributions to calculate the corresponding NLO (LO) unpolarized cross section. We have chosen the scale $`\mu =p_T`$ which leads to the largest cross sections in Fig. 1. We also indicate in Fig. 2 an estimate of the statistical accuracy that might be achievable at RHIC, based on $$\delta A_{\mathrm{TT}}\frac{1}{P_AP_B\sqrt{\sigma _{\mathrm{bin}}}},$$ (12) with beam polarizations $`P_{A,B}`$ of $`70\%`$, and an integrated luminosity $``$ of 320 and $`800\mathrm{pb}^1`$ for c.m.s. energies of $`\sqrt{S}=200`$ and 500 GeV, respectively. $`\sigma _{\mathrm{bin}}`$ denotes the unpolarized cross section integrated over the $`p_T`$-bin for which the error is to be determined. Clearly, the statistics would be sufficient to measure even asymmetries as small as the ones shown in Fig. 2. However, it is likely that the systematic error on spin asymmetries at RHIC will not be much smaller than $`10^3`$, in which case it would appear to be very difficult to access transversity from $`A_{\mathrm{TT}}`$ for single-inclusive pion production. We stress again that the results shown in Fig. 2 are upper bounds, at least within the GRV/GRSV framework with its low input scale for the evolution. If the bound in Eq. (11) turns out to be not saturated at that scale the asymmetries would be even smaller. On the other hand, if we used transversity densities that saturate Eq. (11) at a higher scale, say $`\mu _01\mathrm{GeV}`$, the results for $`A_{\mathrm{TT}}`$ would be somewhat larger. In any case the measurement is very challenging at RHIC. We now turn to transversely polarized $`\overline{p}p`$ collisions with $`\sqrt{S}=14.5`$ GeV at the planned GSI-FAIR facility. We first note that for this rather moderate c.m.s. energy the pion transverse momentum can at most reach $`7.25`$ GeV, at mid rapidity. In our study we integrate over $`1<\eta _{\mathrm{lab}}<2.5`$, where $`\eta _{\mathrm{lab}}`$ is the pseudorapidity of the pion in the laboratory frame. We count positive rapidity in the forward direction of the antiproton. For the asymmetric collider option we consider here, $`\eta _{\mathrm{lab}}`$ is related to the c.m.s. pseudorapidity $`\eta `$ via $$\eta _{\mathrm{lab}}=\eta +\frac{1}{2}\mathrm{ln}\frac{E_{\overline{p}}}{E_p},$$ (13) where $`E_{\overline{p}}`$, $`E_p`$ are the antiproton and proton energies. The rapidity interval we use is roughly symmetric in c.m.s. pseudorapidities, $`\|\eta \|1.75`$. Figure 3 shows our results for the unpolarized (left) and transversely polarized (right) cross sections at NLO and LO, as functions of $`p_T`$. For the calculation in the unpolarized case we have chosen the GRV ref:grv parton distributions. This choice is motivated by our ansatz for the transversity distributions, for which we had also used the GRV densities when saturating the Soffer inequality, Eq. (11). Unlike at RHIC energies, at $`\sqrt{S}=14.5`$ GeV and $`p_T`$ of several GeV, rather large momentum fractions $`x_{a,b}`$ of the partons are probed in Eq. (3), where the polarized and unpolarized parton densities for a given parton type are expected to become similar ref:largex . It then appears most sensible to use the same parton distributions in the unpolarized case that we used when modeling the transversity densities. In this way we avoid any artificial effects in the NLO corrections and $`A_{\mathrm{TT}}`$ induced by a mismatch in the $`x1`$ behavior of the parton densities used in the calculation. The shaded bands in the upper panels of Fig. 3 again indicate the uncertainties due to scale variation in the range $`p_T\mu 4p_T`$. One can see that for both, the unpolarized and the polarized cross sections, the scale dependence does not really improve from LO to NLO. This is a characteristic feature in low-order perturbative calculations of cross sections for lower fixed-target energies, suggesting that corrections beyond NLO are still very significant. Indeed, it was recently shown ref:resumpion that for inclusive-hadron production in the fixed-target regime certain double-logarithmic corrections to the partonic cross sections are important at each order of perturbation theory, and need to be resummed to all orders to achieve an adequate theoretical description. Such a resummation will be required in particular in the case we are considering here and would be very desirable for the future, along with a study of power corrections. We emphasize that when the proposed measurements of $`A_{\mathrm{TT}}`$ will be performed, it will be crucial to have precise measurements also of the unpolarized cross section, in order to test the theoretical framework. Only if the theory is sufficiently understood will data on $`A_{\mathrm{TT}}`$ become useful for determining transversity. Similar comparisons of data ref:rhicdata and theoretical calculations for the unpolarized neutral-pion cross section at RHIC have shown an excellent agreement even down to fairly low pion transverse momenta, which has indeed provided much confidence that the calculations based on partonic hard-scattering are adequate, so that spin asymmetries measured at RHIC determine the spin-dependent parton distributions of the proton. The lower parts of Fig. 3 display the corresponding $`K`$-factors at scale $`\mu =2p_T`$. We note that these decrease as $`p_T`$ increases, which is related entirely to the different behavior of the LO and NLO parton distributions at large $`x`$. Had we chosen the same parton distributions at LO and NLO, the $`K`$-factors would actually slightly increase with $`p_T`$, as a result of the large double-logarithmic corrections mentioned above, which first arise at NLO and are known to enhance the cross section. Figure 4 shows upper bounds–again within the GRV/GRSV framework–for the double-spin asymmetry $`A_{\mathrm{TT}}`$ for the GSI-FAIR facility. The scale $`\mu `$ is again set to $`p_T`$. We also give expectations for the statistical errors that may be achievable in experiment. We have calculated these using Eq. (12), assuming an integrated luminosity of $`=150\mathrm{pb}^1`$, and beam polarizations of $`30\%`$ and $`50\%`$ for the antiprotons and protons, respectively. ## IV Conclusions We have presented in this paper the complete NLO QCD corrections for the partonic hard-scattering cross sections relevant for the double-spin asymmetry $`A_{\mathrm{TT}}`$ for single-inclusive high-$`p_T`$ pion production in collisions of transversely polarized hadrons. This asymmetry could be a tool to determine the transversity distributions of the nucleon. Our calculation is based on a largely analytical evaluation of the NLO partonic cross sections, and we have used a projection technique for treating the characteristic azimuthal-angle dependence introduced by the transverse spin vectors. In our phenomenological studies we found that the spin asymmetry $`A_{\mathrm{TT}}`$ is expected to be very small in $`pp`$ collisions at RHIC and even decreases when going from LO to NLO, due to a larger $`K`$-factor in the unpolarized case. We have also studied $`A_{\mathrm{TT}}`$ for possibly forthcoming transversely polarized $`\overline{p}p`$ collisions in an asymmetric collider mode at the GSI-FAIR facility. Here, the spin asymmetry may be much larger, but it will be crucial in the future to investigate the effects of all-order resummations of large Sudakov logarithms. Detailed measurements of the unpolarized cross sections will be essential for testing the applicability of the theoretical framework at the moderate c.m.s. energies available at GSI-FAIR. ## Acknowledgments We thank F. Rathmann, E. Reya, and W. van Neerven for useful discussions. W.V. is grateful to RIKEN, Brookhaven National Laboratory and the U.S. Department of Energy (contract number DE-AC02-98CH10886) for providing the facilities essential for the completion of this work. This work was supported in part by the “Bundesministerium für Bildung und Forschung (BMBF)” and the “Foundation for Fundamental Research on Matter” (FOM), The Netherlands.
warning/0506/math0506383.html	ar5iv	text	# Parameters Changes for Generalized Power Series ## 1. Introduction The goal of this paper is to provide an algebraic framework for a combinatorial phenomenon arising from a resemblance of variables changes. It is motivated by a well-known formula of Jacobi stated and generalized in this paper as Theorem 4.7 in its modern guise. Our central idea is to find a new notion of differential and a generalization of variable, with which Jacobians appear naturally. To obtain combinatorial information from this algebraic framework, we define an analogue of cohomology residue map . The new residue map also fits our philosophy that a residue comes from a differential. Another recent interpretation of Jacobi’s formula can be found in , where differentials are lacking. To gain a perspective of this paper, it is helpful to look at the method of generating functions whose foundation is built up by rings of formal power series and the operation of equating coefficients. While elegant and easy to implement, the effect of variable changes is not clear in the method without the notion of Kähler differentials. For instance, the role of the Jacobian occurring in a variables change is not transparent. The situation can be improved by meromorphic differentials, with which contour integrations give an alternative way to take coefficients . In the analytic procedure of the method of generating functions, extra attentions are paid to conditions without combinatorial significance such as convergence of sequences and paths for integrations. Removing unnecessary analytic constraints, the author arrives at certain cohomology classes of separated differentials . The process of integration is replaced by cohomology residue maps, which play a significant role in Grothendieck duality theory. Our new algebraic framework is supported by the fact that formal power series rings are in fact a notion free from variables. From this point of view, Lagrange inversion formulae are simply a phenomenon of variables changes ; and pairs of inverse relations are a phenomenon of Schauder bases changes . The cohomology residues come in as an amenable tool for realizing these phenomena. The formalism of cohomology residues is simple. For instance, $$\mathrm{res}\left[\begin{array}{c}\mathrm{\Phi }dX_1\mathrm{}dX_n\\ X_1,\mathrm{},X_n\end{array}\right]=\text{ the constant term of }\mathrm{\Phi }$$ for a power series $`\mathrm{\Phi }\kappa [[X_1,\mathrm{},X_n]]`$ with coefficients in a field $`\kappa `$. Although working well on wide range of problems in combinatorial analysis , a variable change from $`X_i`$ to $`X_i^1`$ is not available. In this paper, we work on a field $`\kappa [[e^𝒢]]`$ of generalized power series with exponents in a totally ordered Abelian group $`𝒢`$ and coefficients in $`\kappa `$ (see Section 2 for a review). The notion of variables is extended to include their inverses. The logarithmic analogue $$\mathrm{res}\left[\begin{array}{c}\mathrm{\Phi }\mathrm{dlog}X_1\mathrm{}\mathrm{dlog}X_n\\ \mathrm{log}X_1,\mathrm{},\mathrm{log}X_n\end{array}\right]$$ of residues is defined, even for a field of positive characteristic. The framework consists of differentials (Section 3), parameters, generalized fractions and residue maps (Section 4). The useful Jacobi’s formula (Theorem 4.7) and Dyson’s conjecture (Section 5) are in fact a phenomenon of parameters changes. These interpretations are seen naturally in our framework. ## 2. Generalized power series We recall the definition and basic properties of generalized power series. For details of proofs, the reader is referred to \[14, Chapter 13, §2\]. Generalized power series are called Malcev-Neumann series in . See for historical remarks on choices of these names. Let $`𝒢`$ be an Abelian group and $`\kappa `$ be a field, whose elements are called scalars. The set $$\kappa \{e^𝒢\}:=\{\underset{g𝒢}{}a_ge^g:a_g\kappa \}$$ of formal sums is a $`\kappa `$-vector space with termwise addition and multiplication (1) $`{\displaystyle \underset{g𝒢}{}}a_ge^g+{\displaystyle \underset{g𝒢}{}}b_ge^g`$ $`:=`$ $`{\displaystyle \underset{g𝒢}{}}(a_g+b_g)e^g,`$ $`b{\displaystyle \underset{g𝒢}{}}a_ge^g`$ $`:=`$ $`{\displaystyle \underset{g𝒢}{}}(ba_g)e^g.`$ For $`\mathrm{\Phi }=_{g𝒢}a_ge^g\kappa \{e^𝒢\}`$, we call $`a_g`$ the $`\kappa `$-coefficient of $`\mathrm{\Phi }`$ at $`e^g`$. The $`\kappa `$-coefficient of $`\mathrm{\Phi }`$ at $`1`$ is called the constant term of $`\mathrm{\Phi }`$ in $`\kappa `$. An element of $`\kappa \{e^𝒢\}`$ is determined by its $`\kappa `$-coefficients. For an element of the form $`Y=e^g`$ with $`g𝒢`$, we use the notation $$\mathrm{log}Y:=g.$$ The support of $`\mathrm{\Phi }`$ is defined as $$\mathrm{supp}\mathrm{\Phi }:=\{g𝒢:\text{ the }\kappa \text{-coefficient of }\mathrm{\Phi }\text{ at }e^g\text{ is not zero}\}.$$ A totally ordered Abelian group is an Abelian group together with a total order compatible with the group structure. ###### Definition 2.1 (generalized power series). Let $`𝒢`$ be a totally ordered Abelian group. We define $$\kappa [[e^𝒢]]:=\{\mathrm{\Phi }\kappa \{e^𝒢\}:\mathrm{supp}\mathrm{\Phi }\text{ is well-ordered}\}.$$ An element in $`\kappa [[e^𝒢]]`$ is called a generalized power series with exponents in $`𝒢`$ and coefficients in $`\kappa `$. Recall that a subset $`A`$ of $`𝒢`$ is well-ordered if every non-empty subset of $`A`$ has a smallest element. ###### Lemma 2.2. Let $`I_1,\mathrm{},I_n`$ be well-ordered subsets of $`𝒢`$ and $`g𝒢`$. The equation $`x_1+\mathrm{}+x_n=g`$ has finitely many solutions $`(x_1,\mathrm{},x_n)`$ with $`x_iI_i`$. The set $`I_1+\mathrm{}+I_n=\{a_1+\mathrm{}+a_n\|a_iI_i\}`$ is well-ordered. So we may define multiplication $$(\underset{g𝒢}{}a_ge^g)(\underset{g𝒢}{}b_ge^g):=\underset{g𝒢}{}\left(\underset{g_1+g_2=g}{}a_{g_1}b_{g_2}\right)e^g$$ for generalized power series. Together with the addition (1), $`\kappa [[e^𝒢]]`$ is a commutative ring with the unit $`1:=e^0`$. A generalized power series $`a_ge^g`$ is positive if $`a_g=0`$ for all $`g0`$. Let $`\mathrm{\Phi }=a_ge^g`$ be a generalized power series. There are no strictly decreasing infinite sequences in $`\mathrm{supp}\mathrm{\Phi }`$. If $`\mathrm{\Phi }`$ is positive, for a fixed $`g𝒢`$, there are only finitely many $`i`$ such that $$\underset{g_1+\mathrm{}+g_i=g}{}a_{g_1}\mathrm{}a_{g_i}0.$$ Given scalars $`c_i`$, we can define a generalized power series $`c_0+c_1\mathrm{\Phi }+c_2\mathrm{\Phi }^2+\mathrm{}`$, whose $`\kappa `$-coefficient at $`e^g`$ is $$\underset{i}{}\left(c_i\underset{g_1+\mathrm{}+g_i=g}{}a_{g_1}\mathrm{}a_{g_i}\right).$$ If the characteristic of $`\kappa `$ is zero, we define $$\mathrm{log}(1+\mathrm{\Phi }):=\underset{\mathrm{}=1}{\overset{\mathrm{}}{}}(1)^{\mathrm{}+1}\frac{\mathrm{\Phi }^{\mathrm{}}}{\mathrm{}}.$$ A non-zero generalized power series $`\mathrm{\Psi }`$ can be factorized uniquely as $$\mathrm{\Psi }=aY(1+\stackrel{~}{\mathrm{\Psi }}),$$ where $`a`$ is a scalar, $`\stackrel{~}{\mathrm{\Psi }}`$ is a positive generalized power series and $`Y=e^g`$ for some $`g𝒢`$. Indeed, $`g`$ is the smallest element of $`\mathrm{supp}\mathrm{\Psi }`$, $`a`$ is the $`\kappa `$-coefficient of $`\mathrm{\Psi }`$ at $`e^g`$ and $`\stackrel{~}{\mathrm{\Psi }}=a^1e^g\mathrm{\Psi }1`$. We call $`a`$ the leading $`\kappa `$-coefficient of $`\mathrm{\Psi }`$. In this paper, the factorization of a non-zero generalized power series refers to the representation of the above form. $`\mathrm{\Psi }`$ is invertible with the inverse $`a^1e^g(1\stackrel{~}{\mathrm{\Psi }}+\stackrel{~}{\mathrm{\Psi }}^2\mathrm{})`$. We conclude that $`\kappa [[e^𝒢]]`$ is a field. ###### Example 2.3 (Laurent series). The field $`\kappa [[e^{}]]`$ with the usual order on $``$ is isomorphic to the field $`\kappa ((X))`$ of Laurent series. ###### Example 2.4 (Hahn , see \[8, p. 445-499\]). The field $`\kappa [[e^{}]]`$ with the usual order on $``$ is of particular interest to algebraic geometers, since it contains an algebraic closure of $`\kappa ((X))`$ if $`\kappa `$ is algebraically closed. ###### Example 2.5 (iterated Laurent series). Let $`𝒢_i=`$ ($`i=1,2`$) and let $`X=e^{(1,0)}`$ and $`Y=e^{(0,1)}`$ be elements of $`\kappa \{e^{}\}`$. With the order $`(m_1,n_1)(m_2,n_2)`$ $``$ $`m_1>m_2\text{ or }m_1=m_2\text{ }\&\text{ }n_1n_2`$ on $`𝒢_1`$, $`\kappa [[e^{𝒢_1}]]`$ is isomorphic to the field $`\kappa ((Y))((X))`$ of iterated Laurent series. With the order $`(m_1,n_1)(m_2,n_2)`$ $``$ $`n_1>n_2\text{ or }n_1=n_2\text{ }\&\text{ }m_1m_2`$ on $`𝒢_2`$, $`\kappa [[e^{𝒢_2}]]`$ is isomorphic to the field $`\kappa ((X))((Y))`$ of iterated Laurent series. As subsets of $`\kappa \{e^{}\}`$, $`\kappa [[e^{𝒢_1}]]`$ and $`\kappa [[e^{𝒢_2}]]`$ are different. The inverse of $`X+Y`$ in $`\kappa [[e^{𝒢_1}]]`$ is $$Y^1(1XY^1+X^2Y^2X^3Y^3+\mathrm{})$$ and that in $`\kappa [[e^{𝒢_2}]]`$ is $$X^1(1X^1Y+X^2Y^2X^3Y^3+\mathrm{}).$$ Let $``$ be a subgroup of $`𝒢`$ with the induced order. In the rest of this paper, we assume that the quotient group $`𝒢/`$ is a free Abelian group of rank $`n`$. In other words, there exist $`u_1,\mathrm{},u_n𝒢`$ such that every element in $`𝒢`$ can be written uniquely as $`h+s_1u_1+\mathrm{}+s_nu_n`$ with $`h`$ and $`s_i`$. We say also that $`𝒢`$ is generated freely by $``$ and $`u_1,\mathrm{},u_n`$. ###### Definition 2.6 (variable). $`e^{u_1},\mathrm{},e^{u_n}`$ are variables of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ if $`𝒢`$ is generated freely by $``$ and $`u_1,\mathrm{},u_n`$. The cardinalities of any sets of variables of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ are the same. We often use the notation $`X_i=e^{u_i}`$. We say also that $`\kappa [[e^𝒢]]`$ is generated by $`\kappa [[e^{}]]`$ and the variables $`X_1,\mathrm{},X_n`$. For an element $$\mathrm{\Psi }=\underset{\stackrel{h}{j_1,\mathrm{},j_n}}{}a_{h,j_1,\mathrm{},j_n}e^hX_1^{j_1}\mathrm{}X_n^{j_n}$$ in $`\kappa [[e^𝒢]]`$ and fixed $`j_1,\mathrm{},j_n`$, we call $`\phi _{j_1,\mathrm{},j_n}:=_ha_{h,j_1,\mathrm{},j_n}e^h`$ the $`\kappa [[e^{}]]`$-coefficient of $`\mathrm{\Psi }`$ at the monomial $`X_1^{j_1}\mathrm{}X_n^{j_n}`$. The $`\kappa [[e^{}]]`$-coefficient of $`\mathrm{\Psi }`$ at $`1`$ (that is, $`\phi _{0,\mathrm{},0}`$) is independent of the choice of variables and is called the constant term of $`\mathrm{\Psi }`$ in $`\kappa [[e^{}]]`$. Indeed, fo any $`h`$, the $`\kappa `$-coefficient of $`\phi _{0,\mathrm{},0}`$ at $`e^h`$ equals to that of $`\mathrm{\Psi }`$ at $`e^h`$. For a given set of variables, $`\mathrm{\Psi }`$ is determined by its $`\kappa [[e^{}]]`$-coefficients. We use the notation $`\mathrm{\Psi }=\mathrm{\Psi }(X_1,\mathrm{},X_n)`$ to indicate that it is represented as the form $$\mathrm{\Psi }=\underset{j_1,\mathrm{},j_n}{}\phi _{j_1,\mathrm{},j_n}X_1^{j_1}\mathrm{}X_n^{j_n}$$ with $`\phi _{j_1,\mathrm{},j_n}\kappa [[e^{}]]`$. ## 3. Differentials A derivation on $`\kappa [[e^𝒢]]`$ is a map $`D`$ from $`\kappa [[e^𝒢]]`$ to a $`\kappa [[e^𝒢]]`$-vector space which satisfies $`D(\mathrm{\Phi }_1+\mathrm{\Phi }_2)=D(\mathrm{\Phi }_1)+D(\mathrm{\Phi }_2)`$ and $`D(\mathrm{\Phi }_1\mathrm{\Phi }_2)=\mathrm{\Phi }_1D(\mathrm{\Phi }_2)+\mathrm{\Phi }_2D(\mathrm{\Phi }_1)`$ for all $`\mathrm{\Phi }_1,\mathrm{\Phi }_2\kappa [[e^𝒢]]`$. Recall that $``$ is a subgroup of $`𝒢`$. A derivation $`D`$ on $`\kappa [[e^𝒢]]`$ is a $`\kappa [[e^{}]]`$-derivation if $`D(\phi )=0`$ for all $`\phi \kappa [[e^{}]]`$. ###### Definition 3.1 (partial derivation). Let $`X_1,\mathrm{},X_n`$ be a set of variables of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$. The partial derivation on $`\kappa [[e^𝒢]]`$ with respect to $`X_i`$ is the well-defined $`\kappa [[e^{}]]`$-derivation $$\frac{}{X_i}:\kappa [[e^𝒢]]\kappa [[e^𝒢]]$$ given by $$\phi _{j_1,\mathrm{},j_n}X_1^{j_1}\mathrm{}X_n^{j_n}j_i\phi _{j_1,\mathrm{},j_n}X_1^{j_1}\mathrm{}X_{i1}^{j_{i1}}X_i^{j_i1}X_{i+1}^{j_{i+1}}\mathrm{}X_n^{j_n}.$$ ###### Definition 3.2 (compatibility with partial derivations). Let $`Y_1,\mathrm{},Y_n`$ be a set of variables of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$. A $`\kappa [[e^{}]]`$-derivation $`D`$ on $`\kappa [[e^𝒢]]`$ is compatible with partial derivations $`/Y_1,\mathrm{},/Y_n`$, if (2) $$D(\mathrm{\Phi })=\underset{i=1}{\overset{n}{}}\frac{\mathrm{\Phi }}{Y_i}D(Y_i)$$ for any $`\mathrm{\Phi }\kappa [[e^𝒢]]`$. $`D`$ is compatible with partial derivations if it is compatible with $`/Y_1,\mathrm{},/Y_n`$ for any set of variables $`Y_1,\mathrm{},Y_n`$. A $`\kappa [[e^{}]]`$-derivation compatible with partial derivations $`/Y_1,\mathrm{},/Y_n`$ is determined by its values at $`Y_1,\mathrm{},Y_n`$. ###### Lemma 3.3. A partial derivation is compatible with partial derivations. ###### Proof. Let $`X_1,\mathrm{},X_n`$ and $`Y_1,\mathrm{},Y_n`$ be two sets of variables of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ with the relation $$\{\begin{array}{cc}Y_i=e^{h_i}X_1^{s_{i1}}\mathrm{}X_n^{s_{in}},\hfill & \\ X_i=e^{h_i^{}}Y_1^{t_{i1}}\mathrm{}Y_n^{t_{in}},\hfill & \end{array}$$ where $`h_i,h_i^{}`$. Note that the matrix $`(t_{ij})`$ is the inverse of $`(s_{ij})`$. To show that $`/X_j`$ is compatible with partial derivations, we check first the relation (2) for $`\mathrm{\Phi }=X_k`$: $`{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{X_k}{Y_i}}{\displaystyle \frac{Y_i}{X_j}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{e^{h_k^{}}(Y_1^{t_{k1}}\mathrm{}Y_n^{t_{kn}})}{Y_i}}{\displaystyle \frac{e^{h_i}(X_1^{s_{i1}}\mathrm{}X_n^{s_{in}})}{X_j}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}t_{ki}s_{ij}{\displaystyle \frac{X_k}{X_j}}=\delta _{kj}={\displaystyle \frac{X_k}{X_j}}.`$ For the general case, it suffices to prove that the coefficients of both sides of (2) at any fixed $`g𝒢`$ are the same. Since the coefficients involve only finitely many $`\kappa [[e^{}]]`$-coefficients at monomials in $`X_1,\mathrm{},X_n`$, the general case is reduced to the special case that $`\mathrm{\Phi }`$ equals to a finite sum of elements of the form $`\phi X_1^{k_1}\mathrm{}X_n^{k_n}`$ with $`\phi \kappa [[e^{}]]`$. From the defining properties of derivations, the special case is further reduced to the case $`\mathrm{\Phi }=X_k`$ that we just proved. ∎ ###### Proposition 3.4 (criterion of compatibility). Let $`D`$ be a $`\kappa [[e^{}]]`$-derivation on $`\kappa [[e^𝒢]]`$. If $`D`$ is compatible with $`/X_1,\mathrm{},/X_n`$ for one set of variables $`X_1,\mathrm{},X_n`$, then $`D`$ is compatible with partial derivations. ###### Proof. Let $`Y_1,\mathrm{},Y_n`$ be another set of variables of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$. The proposition follows from the straightforward computations: $$\underset{i=1}{\overset{n}{}}\frac{\mathrm{\Phi }}{Y_i}D(Y_i)=\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{n}{}}\frac{\mathrm{\Phi }}{Y_i}\frac{Y_i}{X_j}D(X_j)=\underset{j=1}{\overset{n}{}}\frac{\mathrm{\Phi }}{X_j}D(X_j)=D(\mathrm{\Phi }).$$ ###### Definition 3.5 (differentials). A $`\kappa [[e^𝒢]]`$-vector space $`\mathrm{\Omega }_{𝒢/}`$ together with a $`\kappa [[e^{}]]`$-derivation $`d:\kappa [[e^𝒢]]\mathrm{\Omega }_{𝒢/}`$ is the vector space of differentials of $`𝒢`$ over $``$, if $`d`$ is compatible with partial derivations and for any $`\kappa [[e^{}]]`$-derivation $`\delta :\kappa [[e^𝒢]]V`$ compatible with partial derivations, there exists a unique $`\kappa [[e^𝒢]]`$-linear map $`f:\mathrm{\Omega }_{𝒢/}V`$ such that $`fd=\delta `$. In other words, the vector space $`\mathrm{\Omega }_{𝒢/}`$ is the universal object in the category of $`\kappa [[e^{}]]`$-derivations on $`\kappa [[e^𝒢]]`$ compatible with partial derivations. Elements of $`\mathrm{\Omega }_{𝒢/}`$ are called differentials of $`𝒢`$ over $``$, or simply differentials if $`𝒢`$ and $``$ are obvious in the context. They are different from Kähler differentials of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$, which form the universal object $`\mathrm{\Omega }_{\kappa [[e^𝒢]]/\kappa [[e^{}]]}`$ in the category of $`\kappa [[e^{}]]`$-derivations on $`\kappa [[e^𝒢]]`$. For example, we will see in the next proposition that $`\mathrm{\Omega }_{/0}`$ is an one-dimensional $`\kappa [[e^{}]]`$-vector space with the usual order on $``$. However, if the characteristic of $`\kappa `$ is zero, differential basis of $`\mathrm{\Omega }_{\kappa [[e^{}]]/\kappa }`$ (that is, a subset $`B`$ of $`\kappa [[e^{}]]`$ such that $`\{d\mathrm{\Phi }:\mathrm{\Phi }B\}`$ forms a basis of $`\mathrm{\Omega }_{\kappa [[e^{}]]/\kappa }`$) is exactly transcendence basis of $`\kappa ((X))`$ over $`\kappa `$, whose cardinality is infinite. This example shows also that there do exist $`\kappa [[e^{}]]`$-derivations not compatible with partial derivations. ###### Proposition 3.6 (existence of differentials). The vector space of differentials of $`𝒢`$ over $``$ exists (with our assumption that $`𝒢/`$ is free of rank $`n`$). The differentials $`dX_1,\mathrm{},dX_n`$ form a basis of $`\mathrm{\Omega }_{𝒢/}`$ for any set of variables $`X_1,\mathrm{},X_n`$. ###### Proof. Let $`V`$ be a $`\kappa [[e^𝒢]]`$-vector space with basis $`v_1,\mathrm{},v_n`$. The $`\kappa [[e^{}]]`$-derivation $`d:\kappa [[e^𝒢]]V`$ defined by $`d\mathrm{\Phi }=(\mathrm{\Phi }/X_i)v_i`$ is compatible with $`\mathrm{\Phi }/X_1,\mathrm{},\mathrm{\Phi }/X_n`$ and hence compatible with partial derivations. It is easy to see that $`V`$ together with $`d`$ satisfies the universal property. ∎ Let $`\mathrm{\Phi }`$ be a positive generalized power series and $`c_i`$ be scalars. For the special case that there are only finitely many non-zero $`c_i`$, clearly $$d(c_0+c_1\mathrm{\Phi }+c_2\mathrm{\Phi }^2+\mathrm{})=(c_1+2c_2\mathrm{\Phi }+3c_3\mathrm{\Phi }^2+\mathrm{})d\mathrm{\Phi }.$$ For arbitrary $`c_i`$, note that the $`\kappa `$-coefficients of the generalized power series on both sides of the above equation at any $`g𝒢`$ involve only finitely many $`c_i`$. By reduction to the special case, we see that the above equation always holds. In particular $`d\mathrm{log}(1+\mathrm{\Phi })=d\mathrm{\Phi }/(1+\mathrm{\Phi })=d(1+\mathrm{\Phi })/(1+\mathrm{\Phi })`$, if the characteristic of $`\kappa `$ is zero. Even though logarithmic functions are not defined for a field with positive characteristic, we still use the notation $$\mathrm{dlog}\mathrm{\Phi }:=\frac{d\mathrm{\Phi }}{\mathrm{\Phi }}$$ for a non-zero generalized power series $`\mathrm{\Phi }`$ with coefficients in an arbitrary field. The operator $`\mathrm{dlog}`$ transforms the multiplication of non-zero generalized power series to an addition: $$\mathrm{dlog}(\mathrm{\Phi }_1\mathrm{\Phi }_2)=\mathrm{dlog}\mathrm{\Phi }_1+\mathrm{dlog}\mathrm{\Phi }_2.$$ Let $`X_1,\mathrm{},X_n`$ be variables of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$. Given $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n\kappa [[e^𝒢]]`$, we define their Jacobian with respect to $`X_1,\mathrm{},X_n`$ to be $$\left\|\frac{𝚽}{𝐗}\right\|:=\left\|\frac{(\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n)}{(X_1,\mathrm{},X_n)}\right\|:=det\left(\frac{\mathrm{\Phi }_i}{X_j}\right).$$ One is often interested in the $`\kappa [[e^{}]]`$-coefficient of (3) $$\frac{1}{𝚽}\left\|\frac{𝚽}{𝐗}\right\|=\frac{1}{\mathrm{\Phi }_1\mathrm{}\mathrm{\Phi }_n}\left\|\frac{(\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n)}{(X_1,\mathrm{},X_n)}\right\|$$ at $`𝐗^1`$ with the conventions $`𝚽:=\mathrm{\Phi }_1\mathrm{}\mathrm{\Phi }_n`$ and $`𝐗:=X_1\mathrm{}X_n`$. Since a Jacobian appears in the generalized power series, it is more natural to work on the $`n`$th exterior product $`^n\mathrm{\Omega }_{𝒢/}`$ of $`\mathrm{\Omega }_{𝒢/}`$, which is a dimension one $`\kappa [[e^𝒢]]`$-vector space with a basis $$d𝐗:=dX_1\mathrm{}dX_n.$$ For $`\mathrm{\Phi }\kappa [[e^𝒢]]`$ and $`g𝒢`$, the $`\kappa `$-coefficient of $`\mathrm{\Phi }d𝐗`$ at $`e^gd𝐗`$ is defined as the $`\kappa `$-coefficient of $`\mathrm{\Phi }`$ at $`e^g`$; the $`\kappa [[e^{}]]`$-coefficient of $`\mathrm{\Phi }d𝐗`$ at $`X_1^{j_1}\mathrm{}X_n^{j_n}d𝐗`$ is defined as the $`\kappa [[e^{}]]`$-coefficient of $`\mathrm{\Phi }`$ at $`X_1^{j_1}\mathrm{}X_n^{j_n}`$. ###### Proposition 3.7 (vanishing of coefficients). Let $`\mathrm{\Phi }_i,\mathrm{},\mathrm{\Phi }_n`$ be non-zero generalized power series. If some $`i_j`$ is not equal to $`1`$, the $`\kappa [[e^{}]]`$-coefficient of $$\frac{d𝚽}{𝚽^𝐢}:=\frac{d\mathrm{\Phi }_1}{\mathrm{\Phi }_1^{i_1}}\mathrm{}\frac{d\mathrm{\Phi }_n}{\mathrm{\Phi }_n^{i_n}}$$ at $`\mathrm{dlog}𝐗:=d𝐗/𝐗`$ is zero for any set of variables $`X_1,\mathrm{},X_n`$. ###### Proof. The proposition is equivalent to that the $`\kappa `$-coefficient $`c_h`$ of $`d𝚽/𝚽^𝐢`$ at $`e^h\mathrm{dlog}𝐗`$ is zero for any $`h`$. Since $`c_h`$ involves only finitely many non-zero $`\kappa `$-coefficients of $`\mathrm{\Phi }_i`$, we may assume that $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ has only finitely many non-zero $`\kappa `$-coefficients $`a_1,\mathrm{},a_m`$. The coefficient $`c_h`$ is obtained from $`a_1,\mathrm{},a_m`$ by finitely many algebraic operations in $`\kappa `$ (additions, subtractions, multiplications and divisions). There is a polynomial $`f[T_1,\mathrm{},T_m]`$ and $`s`$ such that $$\frac{f(a_1,\mathrm{},a_m)}{(a_1\mathrm{}a_m)^s}=c_h.$$ To show $`f(a_1,\mathrm{},a_m)`$ is zero, it suffices to show that so is $`f(T_1,\mathrm{},T_m)`$. Replacing $`a_i`$ by $`T_i`$, the $`\kappa `$-coefficients of $`\mathrm{\Phi }_i`$ becomes elements in the field $`(T_1,\mathrm{},T_m)`$. So we may assume that $`\kappa =(T_1,\mathrm{},T_m)`$. In particular, $`\kappa `$ has characteristic zero. Following the idea of \[3, Section 1\], we treat first the special case that all $`i_{\mathrm{}}`$ are zero. The derivation $`d`$ is $`\kappa `$-linear, so we may assume furthermore that $`\mathrm{\Phi }_i`$ has only one non-zero $`\kappa `$-coefficient, that is, $`\mathrm{\Phi }_i=a_ie^{h_i}X_1^{s_{i1}}\mathrm{}X_n^{s_{in}}`$ for some $`a_i\kappa `$, $`h_i`$ and $`s_{ij}`$. Under these assumptions, $$d\mathrm{\Phi }_1\mathrm{}d\mathrm{\Phi }_n=a_1\mathrm{}a_n(dets_{ij})e^{{\scriptscriptstyle h_i}}X_1^{1+{\scriptscriptstyle s_{i1}}}\mathrm{}X_n^{1+{\scriptscriptstyle s_{in}}}d𝐗.$$ In order to have non-zero coefficients, $`dets_{ij}`$ can not vanish in $`\kappa `$. But this would imply that the power of some $`X_i`$ in the right hand side of the above equation is not $`1`$. Hence the $`\kappa [[e^{}]]`$-coefficient of $`d𝚽/𝚽^𝐢`$ at $`\mathrm{dlog}𝐗`$ is zero. For the general case, we may assume that $`i_{\mathrm{}}=1`$ for $`\mathrm{}`$ greater than some fixed $`j`$ and $`i_{\mathrm{}}1`$ for $`\mathrm{}j`$. Since the characteristic of $`\kappa `$ is assumed to be zero, $`1i_{\mathrm{}}`$ is invertible in $`\kappa `$ for $`\mathrm{}j`$. The general case is reduced to the special case from the following straightforward computation. $`{\displaystyle \frac{d\mathrm{\Phi }_1}{\mathrm{\Phi }_1^{i_1}}}\mathrm{}{\displaystyle \frac{d\mathrm{\Phi }_n}{\mathrm{\Phi }_n^{i_n}}}`$ $`=`$ $`d\left({\displaystyle \frac{\mathrm{\Phi }_1^{1i_1}}{1i_1}}\right)\mathrm{}d\left({\displaystyle \frac{\mathrm{\Phi }_j^{1i_j}}{1i_j}}\right){\displaystyle \frac{d\mathrm{\Phi }_{j+1}}{\mathrm{\Phi }_{j+1}}}\mathrm{}{\displaystyle \frac{d\mathrm{\Phi }_n}{\mathrm{\Phi }_n}}`$ $`=`$ $`d\left({\displaystyle \frac{\mathrm{\Phi }_1^{1i_1}}{1i_1}}\mathrm{\Phi }_{j+1}^1\mathrm{}\mathrm{\Phi }_n^1\right)d\left({\displaystyle \frac{\mathrm{\Phi }_2^{1i_2}}{1i_2}}\right)\mathrm{}d\left({\displaystyle \frac{\mathrm{\Phi }_j^{1i_j}}{1i_j}}\right)`$ $`d\mathrm{\Phi }_{j+1}\mathrm{}d\mathrm{\Phi }_n.`$ ###### Proposition 3.8 (determinant of exponents). Let $`X_1,\mathrm{},X_n`$ be a set of variables and $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ be non-zero generalized power series with the factorizations $`\mathrm{\Phi }_i=a_ie^{h_i}X_1^{s_{i1}}\mathrm{}X_n^{s_{in}}(1+\stackrel{~}{\mathrm{\Phi }}_i)`$. The $`\kappa [[e^{}]]`$-coefficient of $$\mathrm{dlog}𝚽:=\mathrm{dlog}\mathrm{\Phi }_1\mathrm{}\mathrm{dlog}\mathrm{\Phi }_n=\frac{1}{𝚽}\left\|\frac{𝚽}{𝐗}\right\|d𝐗$$ at $`\mathrm{dlog}𝐗`$ is $`dets_{ij}`$. ###### Proof. As the proof of Proposition 3.7, we may assume that $`\kappa `$ has characteristic zero. The element $`\mathrm{log}(1+\stackrel{~}{\mathrm{\Phi }}_i)`$ can be defined, with which $$\mathrm{dlog}\mathrm{\Phi }_i=\mathrm{dlog}(1+\stackrel{~}{\mathrm{\Phi }}_i)+\underset{j=1}{\overset{n}{}}s_{ij}\frac{dX_j}{X_j}.$$ By Proposition 3.7, the $`\kappa [[e^{}]]`$-coefficient of $`\mathrm{dlog}𝚽`$ at $`\mathrm{dlog}𝐗`$ equals to that of $$\left(\underset{j=1}{\overset{n}{}}s_{1j}\frac{dX_j}{X_j}\right)\mathrm{}\left(\underset{j=1}{\overset{n}{}}s_{nj}\frac{dX_j}{X_j}\right),$$ which is clearly $`dets_{ij}`$. ∎ We would like to define a map independent of the choice of variables with the effect of taking $`\kappa [[e^{}]]`$-coefficients of a generalized power series. Restricting ourselves to $`\kappa [[e^𝒢]]`$ does not work. For instance, if we replace $`X_1`$ by $`X_1/\phi `$ with some non-zero element $`\phi \kappa [[e^{}]]`$, the $`\kappa [[e^{}]]`$-coefficient $`\eta `$ of a generalized power series at $`X_1^{j_1}\mathrm{}X_n^{j_n}`$ becomes $`\eta \phi ^{j_1}`$. Working on $`^n\mathrm{\Omega }_{𝒢/}`$ still has problems: For instance, if we switch the order of two variables, the $`\kappa [[e^{}]]`$-coefficients of a generalized power series change signs. In particular, the $`\kappa [[e^{}]]`$-coefficient of $`dX_1\mathrm{}dX_n`$ at $`dX_2dX_1dX_3\mathrm{}dX_n`$ is $`1`$. In the next section, we will introduce parameters and generalized fractions to achieve our goal. ## 4. Residues Recall that $`𝒢`$ is a totally ordered Abelian group generated freely by a subgroup $``$ and $`n`$ elements. ###### Definition 4.1 (multiplicity). Let $`\mathrm{\Phi }`$ be a non-zero generalized power series with the factorization $`\mathrm{\Phi }=aY(1+\stackrel{~}{\mathrm{\Phi }})`$. The multiplicities of $`\mathrm{\Phi }`$ with respect to a set of variables $`X_1,\mathrm{},X_n`$ are the integers $`i_1,\mathrm{},i_n`$ such that $`Y=e^hX_1^{i_1}\mathrm{}X_n^{i_n}`$ for some $`h`$. ###### Definition 4.2 (parameter). Non-zero generalized power series $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ form a system of parameters (or simply parameters) of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ if the determinant of their multiplicities (with respect to a set of variables) is not zero in $`\kappa `$. The definition is independent of the choice of variables. The determinant of the multiplicities of a system of parameters is not zero in $``$. Let $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ be parameters of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ with the factorizations $`\mathrm{\Phi }_i=a_iY_i(1+\stackrel{~}{\mathrm{\Phi }}_i)`$. The necessary and sufficient condition for $`e^{h_1}Y_1^{i_1}\mathrm{}Y_n^{i_n}=e^{h_2}Y_1^{j_1}\mathrm{}Y_n^{j_n}`$ with $`h_i`$ is $`i_1=j_1`$, $`\mathrm{}`$, $`i_n=j_n`$ and $`h_1=h_2`$. ###### Lemma 4.3. Let $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ be parameters of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$. Then $`\mathrm{dlog}𝚽`$ is a basis of the $`\kappa [[e^𝒢]]`$-vector space $`^n\mathrm{\Omega }_{𝒢/}`$. ###### Proof. Since $`^n\mathrm{\Omega }_{𝒢/}`$ has dimension one, we need to check that $`\mathrm{dlog}𝚽0`$. Let $`X_1,\mathrm{},X_n`$ be variables. Replacing $`X_i`$ by its inverse if $`\mathrm{log}X_i<0`$, we may assume that $`\mathrm{log}X_i>0`$ for all $`i`$. With the factorizations $`\mathrm{\Phi }_i=a_ie^{h_i}X_1^{s_{i1}}\mathrm{}X_n^{s_{in}}(1+\stackrel{~}{\mathrm{\Phi }}_i)`$, we can write $`d\mathrm{log}𝚽`$ $`=`$ $`({\displaystyle \frac{d\stackrel{~}{\mathrm{\Phi }}_1}{1+\stackrel{~}{\mathrm{\Phi }}_1}}+{\displaystyle \underset{j}{}}s_{1j}{\displaystyle \frac{dX_j}{X_j}})\mathrm{}({\displaystyle \frac{d\stackrel{~}{\mathrm{\Phi }}_n}{1+\stackrel{~}{\mathrm{\Phi }}_n}}+{\displaystyle \underset{j}{}}s_{nj}{\displaystyle \frac{dX_j}{X_j}})`$ $`=`$ $`\mathrm{\Psi }d𝐗+{\displaystyle \frac{dets_{ij}}{𝐗}}d𝐗`$ for some $`\mathrm{\Psi }\kappa [[e^𝒢]]`$. Note that, from our convention of positivity of $`\stackrel{~}{\mathrm{\Phi }}_i`$, the support of $`\mathrm{\Psi }`$ consists of only elements greater than $`\mathrm{log}𝐗`$. Since $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ are parameters, the leading coefficient $`dets_{ij}`$ of $`\mathrm{\Psi }+(dets_{ij})𝐗^1`$ is not zero. Therefore $`\mathrm{dlog}𝚽0`$. ∎ Let $`V`$ be a $`\kappa [[e^𝒢]]`$-vector space. In the set $$\{(\alpha ,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n)V\times \kappa [[e^𝒢]]^n:\text{ }\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n\text{ are parameters}\},$$ we define an equivalence relation: $$(\alpha ,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n)(\beta ,\mathrm{\Psi }_1,\mathrm{},\mathrm{\Psi }_n)\frac{\beta }{dett_{ij}}=detu_{ij}\frac{\alpha }{dets_{ij}},$$ where $`s_{ij}`$ (resp. $`t_{ij}`$) are multiplicities of $`\mathrm{\Phi }_i`$ (resp. $`\mathrm{\Psi }_i`$) with respect to a set of variables $`X_1,\mathrm{},X_n`$ (resp. $`Y_1,\mathrm{},Y_n`$) and $`Y_i=e^{h_i}X_1^{u_{i1}}\mathrm{}X_n^{u_{in}}`$. The equivalence relation is independent of the choices of variables. ###### Definition 4.4 (generalized fraction). A generalized fraction $$\left[\begin{array}{c}\alpha \\ \mathrm{log}𝚽\end{array}\right]:=\left[\begin{array}{c}\alpha \\ \mathrm{log}\mathrm{\Phi }_1,\mathrm{},\mathrm{log}\mathrm{\Phi }_n\end{array}\right]$$ is the equivalence class containing $`(\alpha ,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n)`$. We call $`\alpha `$ the numerator of the generalized fraction. The set of generalized fractions with numerators in $`V`$ is denoted by $`\mathrm{H}(V)`$. We choose the notation $`\mathrm{H}(V)`$, because it might relate to some cohomology object as the case of the theory of local cohomology residues for formal power series rings. ###### Definition 4.5 (residue). Let $`X_1,\mathrm{},X_n`$ be variables of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$. We define the residue map $$\mathrm{res}_{X_1,\mathrm{},X_n}:\mathrm{H}(^n\mathrm{\Omega }_{𝒢/})\kappa [[e^{}]]$$ with respect to $`X_1,\mathrm{},X_n`$ by $$\mathrm{res}_{X_1,\mathrm{},X_n}\left[\begin{array}{c}\mathrm{\Phi }\mathrm{dlog}𝐗\\ \mathrm{log}𝐗\end{array}\right]=\text{ the constant term of }\mathrm{\Phi }\text{ in }\kappa [[e^{}]]\text{,}$$ where $`\mathrm{\Phi }\kappa [[e^𝒢]]`$. ###### Proposition 4.6 (invariance of residues). $`\mathrm{res}_{X_1,\mathrm{},X_n}=\mathrm{res}_{Y_1,\mathrm{},Y_n}`$ for any two sets of variables $`X_1,\mathrm{},X_n`$ and $`Y_1,\mathrm{},Y_n`$ of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$. ###### Proof. Write $`Y_i=e^{h_i}X_1^{s_{i1}}\mathrm{}X_n^{s_{in}}`$. Then $$\mathrm{dlog}𝐘=\mathrm{dlog}(X_1^{s_{11}}\mathrm{}X_n^{s_{1n}})\mathrm{}\mathrm{dlog}(X_1^{s_{n1}}\mathrm{}X_n^{s_{nn}})=(dets_{ij})\mathrm{dlog}𝐗.$$ For any $`\mathrm{\Phi }\kappa [[e^𝒢]]`$, $$\left[\begin{array}{c}\mathrm{\Phi }\mathrm{dlog}𝐘\\ \mathrm{log}𝐘\end{array}\right]=\left[\begin{array}{c}(dets_{ij})^1\mathrm{\Phi }\mathrm{dlog}𝐘\\ \mathrm{log}𝐗\end{array}\right]=\left[\begin{array}{c}\mathrm{\Phi }\mathrm{dlog}𝐗\\ \mathrm{log}𝐗\end{array}\right].$$ Therefore $`\mathrm{res}_{X_1,\mathrm{},X_n}=\mathrm{res}_{Y_1,\mathrm{},Y_n}`$. ∎ Taking the residue is a map equating coefficients independent of the choice of variables. We denote $`\mathrm{res}:=\mathrm{res}_{X_1,\mathrm{},X_n}`$. Let $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n\kappa [[e^𝒢]]`$ and $`\phi _{i_1,\mathrm{},i_n}\kappa [[e^{}]]`$, where the indices $`i_{\mathrm{}}`$. We assume that there are only finitely many $`\phi _{i_1,\mathrm{},i_n}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}`$ whose support contains any fixed $`g𝒢`$. Under such an assumption, we can define an element $$\mathrm{\Psi }(𝚽):=\phi _{i_1,\mathrm{},i_n}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}\kappa \{e^𝒢\},$$ whose $`\kappa `$-coefficient at $`e^g`$ is the sum of the $`\kappa `$-coefficients of all $`\phi _{i_1,\mathrm{},i_n}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}`$ at $`e^g`$. We say that $`\mathrm{\Psi }(𝚽)`$ is represented by $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ with $`\kappa [[e^{}]]`$-coefficient $`\phi _{i_1,\mathrm{},i_n}`$ at $`\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}`$. Given $`\phi _{i_1,\mathrm{},i_n}`$, whether or not $`\mathrm{\Psi }(𝚽)\kappa [[e^𝒢]]`$ depends on the sequence $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$. For instance, as seen in Example 2.5, the inverse of $`X+Y`$ in the iterated Laurent series $`\kappa [[e^{𝒢_1}]]`$ is $$\mathrm{\Psi }(X,Y):=Y^1XY^2+X^2Y^3\mathrm{}.$$ However, $$\mathrm{\Psi }(Y,X)=X^1YX^2+Y^2X^3\mathrm{}$$ although defined is not contained in $`\kappa [[e^{𝒢_1}]]`$. If there are only finitely many nonzero $`\kappa [[e^{}]]`$-coefficients for a representation $`\mathrm{\Psi }(𝚽)`$ of an element in $`\kappa \{e^𝒢\}`$ and the indices of these nonzero coefficients are all non-negative, we say that $`\mathrm{\Psi }(𝚽)`$ is a polynomial in $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ with coefficients in $`\kappa [[e^{}]]`$. The set of these polynomials is a subring of $`\kappa [[e^𝒢]]`$ denoted by $`\kappa [[e^{}]][\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n]`$, which is exactly the image of the homomorphism $$\kappa [[e^{}]][Y_1,\mathrm{},Y_n]\kappa [[e^𝒢]]$$ of $`\kappa [[e^{}]]`$-algebras sending $`Y_i`$ to $`\mathrm{\Phi }_i`$. For $`F\kappa [[e^{}]][Y_1,\mathrm{},Y_n]`$ sending to $`\mathrm{\Psi }(𝚽)`$ under this homomorphism, we use also the notation $`F(𝚽):=\mathrm{\Psi }(𝚽)`$. If we assume furthermore that $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ are parameters, then the above homomorphism is ono-to-one. A representation $`\mathrm{\Psi }(𝚽)`$ of an element in $`\kappa [[e^𝒢]]`$ is a rational function in $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ with coefficients in $`\kappa [[e^{}]]`$ if there exist $`F_1,F_2\kappa [[e^{}]][Y_1,\mathrm{},Y_n]`$ with $`F_10`$ such that $`F_1(𝚽)\mathrm{\Psi }(𝚽)=F_2(𝚽)`$. Let $`F=F_1/F_2\kappa [[e^{}]](Y_1,\mathrm{},Y_n)`$. We use also the notation $`F(𝚽):=\mathrm{\Psi }(𝚽)`$. Now we interpret and generalize Jacobi’s formula. ###### Theorem 4.7 (Jacobi). Given a representation $`\mathrm{\Psi }(𝚽)=\phi _{i_1,\mathrm{},i_n}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}\kappa [[e^𝒢]]`$ by parameters $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ with $`\phi _{i_1,\mathrm{},i_n}\kappa [[e^{}]]`$, $$\mathrm{res}\left[\begin{array}{c}\mathrm{\Psi }(𝚽)\mathrm{dlog}𝚽\\ \mathrm{log}𝚽\end{array}\right]=\phi _{0,\mathrm{},0}.$$ ###### Proof. Let $`s_{ij}`$ be the multiplicities of $`𝚽`$ with respect to a set of variables $`𝐗`$. Since $$\left[\begin{array}{c}\mathrm{\Psi }(𝚽)\mathrm{dlog}𝚽\\ \mathrm{log}𝚽\end{array}\right]=\left[\begin{array}{c}\frac{\mathrm{\Psi }(𝚽)}{𝚽}\left\|\frac{𝚽}{𝐗}\right\|d𝐗\\ \mathrm{log}𝚽\end{array}\right]=\left[\begin{array}{c}\frac{𝐗}{dets_{ij}}\frac{\mathrm{\Psi }(𝚽)}{𝚽}\left\|\frac{𝚽}{𝐗}\right\|\mathrm{dlog}𝐗\\ \mathrm{log}𝐗\end{array}\right],$$ the theorem is equivalent to the claim that the $`\kappa `$-coefficients of $`(dets_{ij})\phi _{0,\mathrm{},0}`$ at $`e^h`$ and $`(\mathrm{\Psi }(𝚽)/𝚽)\|𝚽/𝐗\|`$ at $`e^h𝐗^1`$ are the same for any $`h`$. The latter involves only finitely many $`\phi _{i_1,\mathrm{},i_n}`$, so we may assume that only finitely many $`\phi _{i_1,\mathrm{},i_n}`$ are not zero. From linearity, we may assume furthermore that $`\mathrm{\Psi }=X_1^{i_1}\mathrm{}X_n^{i_n}`$. The theorem in such a special case was proved in Propositions 3.7 and 3.8. ∎ Another interpretation and generalization of Jacobi’s formula in characteristic zero can be found in \[16, Theorem 3.7\]. While investigates the interplay of two fields, we work on one vector space of differentials. In our approach, combinatorial information appears naturally through a residue map with Jacobians resulted from parameters changes. In , Jacobi’s formula is called a residue theorem. However, residue theorem usually refers to Cauchy’s theorem, which counts residues of a meromorphic function in a region. As a global result relating the poles of a meromorphic function, Cauchy’s residue theorem is considered in a very general context by Grothendieck in algebraic geometry. Jacobi’s formula, exploring parameters changes of one point, is merely a local property! Formulae of the Lagrange inversion type can be studied in the field of generalized power series. Along this direction, one needs to know whether or not every generalized power series can be represented by a system of parameters and in what sense a representation is unique. Let $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ be a system of parameters of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ with the factorizations $`\mathrm{\Phi }_i=a_iY_i(1+\stackrel{~}{\mathrm{\Phi }}_i)`$. The following uniqueness property is obvious: If $$\phi _{i_1,\mathrm{},i_n}^{(1)}Y_1^{i_1}\mathrm{}Y_n^{i_n}=\phi _{i_1,\mathrm{},i_n}^{(2)}Y_1^{i_1}\mathrm{}Y_n^{i_n}\kappa [[e^𝒢]],$$ then $`\phi _{i_1,\mathrm{},i_n}^{(1)}=\phi _{i_1,\mathrm{},i_n}^{(2)}`$ for all $`i_1,\mathrm{},i_n`$. ###### Definition 4.8 (regular parameter). A system of parameters $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ with the factorizations $`\mathrm{\Phi }_i=a_iY_i(1+\stackrel{~}{\mathrm{\Phi }}_i)`$ is regular if, for every element $`\mathrm{\Psi }\kappa [[e^𝒢]]`$, there exists an unique element $`\phi _{i_1,\mathrm{},i_n}Y_1^{i_1}\mathrm{}Y_n^{i_n}\kappa [[e^𝒢]]`$ such that $`\mathrm{\Psi }=\phi _{i_1,\mathrm{},i_n}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}`$. Clearly, variables are regular parameters. ###### Proposition 4.9 (characterization of regularity). A system of parameters of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ is regular if and only if the determinant of their multiplicities (with respect to a set of variables) is invertible in $``$. ###### Proof. Let $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ be parameters of $`\kappa [[e^𝒢]]`$ over $`\kappa [[e^{}]]`$ with the factorizations $`\mathrm{\Phi }_i=a_iY_i(1+\stackrel{~}{\mathrm{\Phi }}_i)`$. Assume that the determinant of their multiplicities is invertible in $``$. This assumption is equivalent to that $`Y_1,\mathrm{},Y_n`$ are variables. We need to find the $`\kappa `$-coefficient $`a_g`$ of a generalized power series $`\phi _{i_1,\mathrm{},i_n}Y_1^{i_1}\mathrm{}Y_n^{i_n}`$ at $`e^g`$ for each $`g𝒢`$ to represent a given generalized power series $`\mathrm{\Psi }=b_ge^g`$. Let $$A=\mathrm{supp}\stackrel{~}{\mathrm{\Phi }}_1+\mathrm{}+\mathrm{supp}\stackrel{~}{\mathrm{\Phi }}_n,$$ $`iA`$ be the sum of $`i`$ copies of $`A`$ for $`i>0`$, and $`0A:=\{0\}`$. The well-ordered set $`\overline{A}:=_{i0}iA`$ contains $`\mathrm{supp}(1+\stackrel{~}{\mathrm{\Phi }}_1)^{i_1}\mathrm{}(1+\stackrel{~}{\mathrm{\Phi }}_n)^{i_n}`$. Let $`B`$ be a well-ordered set containing $`\mathrm{supp}\mathrm{\Psi }`$. If $`g\overline{A}+B`$, we define $`a_g=0`$. For $`g\overline{A}+B`$, we consider the equation $`x+y=g`$ subject to the constraints $`x\overline{A}+B`$ and $`y\mathrm{supp}(1+\stackrel{~}{\mathrm{\Phi }}_1)^{i_1}\mathrm{}(1+\stackrel{~}{\mathrm{\Phi }}_n)^{i_n}`$, where $`i_1,\mathrm{},i_n`$ are integers satisfy $`x=h+i_1\mathrm{log}Y_1+\mathrm{}+i_n\mathrm{log}Y_n`$ for some $`h`$. By Lemma 2.2, the equation with the constraints has finitely many solutions. If $`(x,y)=(g,0)`$ is the only solution, for instance if $`g`$ is the smallest element of $`\overline{A}+B`$, we define $`a_g:=b_g`$. If it has other solutions, say $`(g_{11},g_{21})`$, … , $`(g_{1m},g_{2m})`$ besides $`(g,0)`$, we would like to define $$a_g:=b_g\underset{\mathrm{}=1}{\overset{m}{}}a_{g_1\mathrm{}}c_{\mathrm{}}$$ inductively in terms of $`a_{g_1\mathrm{}}`$, where $`g_1\mathrm{}=h_{\mathrm{}}+i_\mathrm{}1\mathrm{log}Y_1+\mathrm{}+i_\mathrm{}n\mathrm{log}Y_n`$, $`h_{\mathrm{}}`$, $`i_\mathrm{}1,\mathrm{},i_\mathrm{}n`$ and $`c_{\mathrm{}}`$ is the $`\kappa `$-coefficient of $`(1+\stackrel{~}{\mathrm{\Phi }}_1)^{i_\mathrm{}1}\mathrm{}(1+\stackrel{~}{\mathrm{\Phi }}_n)^{i_\mathrm{}n}`$ at $`e^{g_2\mathrm{}}`$. To see the inductive process working, we observe that $`g_1\mathrm{}<g`$, since $`g_2\mathrm{}>0`$. Moreover, if the equation $`x+y=g_1\mathrm{}^{(1)}`$ for each $`g_1\mathrm{}^{(1)}:=g_1\mathrm{}`$ subject to the constraints above has only one solution, $`a_{g_1\mathrm{}}`$ has been defined. Let $`(g_{11}^{(2)},g_{21}^{(2)})`$, $`(g_{12}^{(2)},g_{22}^{(2)})`$, $`(g_{13}^{(2)},g_{23}^{(2)})`$ … be solutions of other equations if any. We repeat the process for the equations $`x+y=g_1\mathrm{}^{(2)}`$ with the same constraints. If these equations have more than one solution, we continue the process. The process has to stopped in finitely many steps, since the elements $`g_1\mathrm{}^{(i)}`$ obtained are contained in $`\overline{A}+B`$, which consists of no strictly decreasing infinite sequences. Therefore $`a_g`$ is defined. From the construction, $`\mathrm{\Psi }`$ is represented by $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ with $`\phi _{i_1,\mathrm{},i_n}`$, where $`a_ge^g=\phi _{i_1,\mathrm{},i_n}Y_1^{i_1}\mathrm{}Y_n^{i_n}\kappa [[e^𝒢]]`$. Assume that there are two representations $$\mathrm{\Psi }=\phi _{i_1,\mathrm{},i_n}^{(1)}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}=\phi _{i_1,\mathrm{},i_n}^{(2)}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n},$$ with $`\mathrm{\Psi }_i=\phi _{i_1,\mathrm{},i_n}^{(i)}Y_1^{i_1}\mathrm{}Y_n^{i_n}\kappa [[e^𝒢]]`$. In the above process, we may take $`B=\mathrm{supp}\mathrm{\Psi }\mathrm{supp}\mathrm{\Psi }_1\mathrm{supp}\mathrm{\Psi }_2`$, which contains both $`\mathrm{supp}\mathrm{\Psi }_1`$ and $`\mathrm{supp}\mathrm{\Psi }_2`$. As $`a_g`$ is determined by $`b_g`$ for $`g\overline{A}+B`$, the representations must be the same. Now we assume that $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ are regular parameters. Let $`X_1,\mathrm{},X_n`$ be a set of variables. There exist $`\mathrm{\Psi }_{\mathrm{}}=\phi _{i_1,\mathrm{},i_n}^{(\mathrm{})}Y_1^{i_1}\mathrm{}Y_n^{i_n}\kappa [[e^𝒢]]`$ such that $`X_{\mathrm{}}=\phi _{i_1,\mathrm{},i_n}^{(\mathrm{})}\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}`$. Since $`Y_1,\mathrm{},Y_n`$ are parameters, the minimal element of $`\mathrm{supp}\mathrm{\Psi }_{\mathrm{}}`$ is $`\mathrm{log}X_{\mathrm{}}`$. This implies that $`X_{\mathrm{}}=e^h_{\mathrm{}}Y_1^{i_\mathrm{}1}\mathrm{}Y_n^{i_\mathrm{}n}`$ for some $`h_{\mathrm{}}`$ and $`i_\mathrm{}1,\mathrm{},i_\mathrm{}n`$. Therefore the determinant of the multiplicities of $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ is invertible in $``$. ∎ The theme of Lagrange inversions in the context of generalized power series is the interrelations between two systems of regular parameters. Let $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ be a system of regular parameters represented by another system of regular parameters $`\mathrm{{\rm Y}}_1,\mathrm{},\mathrm{{\rm Y}}_n`$. The expression $$\mathrm{res}\left[\begin{array}{c}\frac{\mathrm{{\rm Y}}_{\mathrm{}}}{\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}}\mathrm{dlog}𝚽\\ \mathrm{log}𝚽\end{array}\right]$$ gives the $`\kappa [[e^{}]]`$-coefficient of $`\mathrm{{\rm Y}}_{\mathrm{}}`$ at $`\mathrm{\Phi }_1^{i_1}\mathrm{}\mathrm{\Phi }_n^{i_n}`$. Properties of generalized fraction and residues can be used to compute the coefficient in terms of $`\kappa [[e^{}]]`$-coefficients of $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$ at monomials in $`\mathrm{{\rm Y}}_1,\mathrm{},\mathrm{{\rm Y}}_n`$. ## 5. Dyson’s conjecture Let $`a_1,\mathrm{},a_n`$ be non-negative integers. Dyson’s conjecture that $$\text{the constant term of }\underset{1ijn}{}(1\frac{X_i}{X_j})^{a_i}=\frac{(a_1+\mathrm{}+a_n)!}{a_1!\mathrm{}a_n!}$$ was confirmed by Wilson and Gunson independently of each other. We interpret two known proofs of the Dyson’s conjecture in terms of generalized power series with coefficients in $``$ and exponents in $`^n`$, which has a total order compatible with the group structure (for instance, the lexicographic order). Let $`X_1,\mathrm{},X_n`$ be variables of $`[[e^^n]]`$ over $`=[[e^0]]`$. In the first proof, we assume that the variables satisfy $`\mathrm{log}X_1>\mathrm{}>\mathrm{log}X_n`$. Let $`\mathrm{\Phi }_i=_{j=1,ji}^n(1X_i/X_j)^1`$. Using Lagrange interpolation, one can show $`_{i=1}^n\mathrm{\Phi }_i=1`$. Wilson’s proof to the Dyson’s conjecture is based on the parameters $`X_1,\mathrm{\Phi }_2,\mathrm{},\mathrm{\Phi }_n`$, whose multiplicities with respect to $`X_1,\mathrm{},X_n`$ have determinant $$det\left(\begin{array}{ccccc}1& 0& 0& \mathrm{}& 0\\ 1& 1& 0& \mathrm{}& 0\\ 1& 1& 2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 1& 1& 1& \mathrm{}& (n1)\end{array}\right)=(n1)!(1)^{n1}.$$ Wilson’s computation \[15, Proof of Lemma 3\] carried over to our context shows $$\mathrm{dlog}X_1\mathrm{dlog}\mathrm{\Phi }_2\mathrm{}\mathrm{dlog}\mathrm{\Phi }_n=c(n1)!(1)^{n1}\mathrm{\Phi }_1\mathrm{dlog}𝐗$$ for some $`c\kappa `$. The scalar $`c`$ is not zero, since $`\mathrm{dlog}X_1\mathrm{dlog}\mathrm{\Phi }_2\mathrm{}\mathrm{dlog}\mathrm{\Phi }_n`$ generates $`^n\mathrm{\Omega }_{𝒢/}`$. Let $`\mathrm{\Psi }(X_1,\mathrm{},X_n)=X_2^{a_2}\mathrm{}X_n^{a_n}(1X_2\mathrm{}X_n)^{a_1}`$. What we need to compute is the constant term of $`\mathrm{\Psi }(X_1,\mathrm{\Phi }_2,\mathrm{},\mathrm{\Phi }_n)`$, that is, the residue of $`\left[\begin{array}{c}\mathrm{\Psi }(X_1,\mathrm{\Phi }_2,\mathrm{},\mathrm{\Phi }_n)\mathrm{dlog}𝐗\\ \mathrm{log}𝐗\end{array}\right]`$ $`=`$ $`\left[\begin{array}{c}\frac{c^1\mathrm{\Psi }(X_1,\mathrm{\Phi }_2,\mathrm{},\mathrm{\Phi }_n)}{1\mathrm{\Phi }_2\mathrm{}\mathrm{\Phi }_n}\mathrm{dlog}X_1\mathrm{dlog}\mathrm{\Phi }_2\mathrm{}\mathrm{dlog}\mathrm{\Phi }_n\\ \mathrm{log}X_1,\mathrm{log}\mathrm{\Phi }_2,\mathrm{},\mathrm{log}\mathrm{\Phi }_n\end{array}\right].`$ By Theorem 4.7, the constant term of $`\mathrm{\Psi }(X_1,\mathrm{\Phi }_2,\mathrm{},\mathrm{\Phi }_n)`$ is the same as that of $$\frac{c^1\mathrm{\Psi }(X_1,\mathrm{},X_n)}{1X_2\mathrm{}X_n}=\frac{c^1}{X_2^{a_2}\mathrm{}X_n^{a_n}}\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{k+a_1}{a_1}\right)(X_2+\mathrm{}+X_n)^k,$$ which is $$c^1\left(\genfrac{}{}{0pt}{}{a_1+\mathrm{}+a_n}{a_1}\right)\left(\genfrac{}{}{0pt}{}{a_2+\mathrm{}+a_n}{a_2,\mathrm{},a_n}\right)=c^1\frac{(a_1+\mathrm{}+a_n)!}{a_1!\mathrm{}a_n!}$$ occurring when $`k=a_2+\mathrm{}+a_n`$. Now Dyson’s conjecture for the trivial case $`a_0=\mathrm{}=a_n=0`$ shows $`c=1`$. In the second proof, we assume that $`\mathrm{log}X_1<\mathrm{}<\mathrm{log}X_n`$. Following Egorychev , we use the parameters $`\mathrm{{\rm Y}}_i=(1)^{i1}X_i^{n1}_{j<k,ji,ki}(X_jX_k)`$, whose multiplicities with respect to $`X_1,\mathrm{},X_n`$ have determinant $$det\left(\begin{array}{ccccc}n1& n2& n3& \mathrm{}& 0\\ n2& n1& n3& \mathrm{}& 0\\ n2& n3& n1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ n2& n3& n4& \mathrm{}& n1\end{array}\right)=\frac{n!(n1)}{2}$$ \[16, proof of Theorem 5.3\]. (The matrix has diagonal entries $`n1`$ and other entries in each row, except the diagonal, are $`n2,n3,\mathrm{},0`$ from left to right.) By Cramer’s rule, $$\frac{\mathrm{{\rm Y}}_1}{X_1^i}+\mathrm{}+\frac{\mathrm{{\rm Y}}_n}{X_n^i}=\{\begin{array}{cc}\mathrm{\Delta }:=_{j<k}(X_jX_k),\hfill & \text{ if }i=0\text{;}\hfill \\ 0,\hfill & \text{ if }i=1,\mathrm{},n1\text{.}\hfill \end{array}$$ Applying the derivation, we obtain $`\mathrm{{\rm Y}}_1\mathrm{dlog}\mathrm{{\rm Y}}_1+\mathrm{}+\mathrm{{\rm Y}}_n\mathrm{dlog}\mathrm{{\rm Y}}_n`$ $`=`$ $`X_1{\displaystyle \frac{\mathrm{\Delta }}{X_1}}\mathrm{dlog}X_1+\mathrm{}+X_n{\displaystyle \frac{\mathrm{\Delta }}{X_n}}\mathrm{dlog}X_n,`$ $`{\displaystyle \frac{\mathrm{{\rm Y}}_1}{X_1^i}}\mathrm{dlog}\mathrm{{\rm Y}}_1+\mathrm{}+{\displaystyle \frac{\mathrm{{\rm Y}}_n}{X_n^i}}\mathrm{dlog}\mathrm{{\rm Y}}_n`$ $`=`$ $`i{\displaystyle \frac{\mathrm{{\rm Y}}_1}{X_1^i}}\mathrm{dlog}X_1+\mathrm{}+i{\displaystyle \frac{\mathrm{{\rm Y}}_n}{X_n^i}}\mathrm{dlog}X_n`$ ($`i=1,\mathrm{},n1`$). Exterior products of the above elements give rise to $$\frac{\mathrm{{\rm Y}}_1\mathrm{}\mathrm{{\rm Y}}_n\mathrm{\Delta }}{(X_1\mathrm{}X_n)^{n1}}\mathrm{dlog}𝚼=\frac{(n1)!\mathrm{{\rm Y}}_1\mathrm{}\mathrm{{\rm Y}}_n}{(X_1\mathrm{}X_n)^{n1}}(X_1\frac{\mathrm{\Delta }}{X_1}+\mathrm{}+X_n\frac{\mathrm{\Delta }}{X_n})\mathrm{dlog}𝐗.$$ Since $$X_1\frac{\mathrm{\Delta }}{X_1}+\mathrm{}+X_n\frac{\mathrm{\Delta }}{X_n}=\left(\genfrac{}{}{0pt}{}{n}{2}\right)\mathrm{\Delta },$$ the combinatorial number $`n!(n1)/2`$ is exactly compensated in the identity $$\mathrm{dlog}𝚼=\frac{n!(n1)}{2}\mathrm{dlog}𝐗.$$ As , we condiser $`\mathrm{\Psi }(X_1,\mathrm{},X_n)=(X_1+\mathrm{}+X_n)^{a_1+\mathrm{}+a_n}/(X_1^{a_1}\mathrm{}X_n^{a_n})`$. Since $$\underset{i=1,ij}{\overset{n}{}}(1\frac{X_i}{X_j})=\frac{\mathrm{{\rm Y}}_1+\mathrm{}+\mathrm{{\rm Y}}_n}{\mathrm{{\rm Y}}_j},$$ we need to show $$\mathrm{res}\left[\begin{array}{c}\mathrm{\Psi }(𝚼)\mathrm{dlog}𝐗\\ \mathrm{log}𝐗\end{array}\right]=\frac{(a_1+\mathrm{}+a_n)!}{a_1!\mathrm{}a_n!}.$$ This is a special case of Theorem 4.7, since $$\left[\begin{array}{c}\mathrm{\Psi }(𝚼)\mathrm{dlog}𝐗\\ \mathrm{log}𝐗\end{array}\right]=\left[\begin{array}{c}\mathrm{\Psi }(𝚼)\mathrm{dlog}𝚼\\ \mathrm{log}𝚼\end{array}\right]$$ and the constant term of $`\mathrm{\Psi }(X_1,\mathrm{},X_n)`$ is $`(a_1+\mathrm{}+a_n)!/(a_1!\mathrm{}a_n!)`$. One more proof by local cohomology residues is available. See \[11, Identity 14\].
warning/0506/astro-ph0506121.html	ar5iv	text	# Dynamical evolution of neutrino–cooled accretion disks: detailed microphysics, lepton-driven convection, and global energetics ## 1 Introduction Gas accretion by a concentration of mass is an efficient way to transform gravitational binding energy into radiation (Salpeter, 1964; Zel’dovich, 1964), and is responsible for observable phenomena at all scales in astrophysics. The more compact the object, the greater the efficiency. In many cases, systems are observed in steady or quasi–steady state during this process (e.g, CVs, LMXBs and AGNs). Some external agent (e.g., the interstellar medium, a companion star) provides a continuous supply of mass, energy and angular momentum over a timescale that is much longer than the accretion time. The energy dissipated by the flow before it is accreted, and more importantly, what happens to this energy, determines many of the properties of the flow itself. In most cases cooling is present through some radiative electromagnetic process, which acts as a sink, removing the dissipated energy. In the standard thin–disk theory developed by Shakura & Sunyaev (1973), it is extremely efficient, maintaining a “cool” (in the sense that $`kTGMm_p/r`$, where $`M`$ is the mass of the central object) and thin (with a scale height $`Hr`$), nearly Keplerian disk. The requirements of steady state, hydrostatic equilibrium in the vertical direction, energy and angular momentum balance, and an equation of state yield a solution for all relevant variables as a function of the disk radius, $`r`$. The key question of angular momentum transport was addressed by Shakura & Sunyaev with their famous $`\alpha `$ prescription, which allows a parametrization of the viscous stresses and energy dissipation rates. Magnetic fields could be at the physical origin this effect, through the magneto–rotational instability (MRI, Balbus & Hawley, 1991; Hawley & Balbus, 1991). In the case of inefficient cooling, a solution was found in the 1970s for so–called slim disks (Shapiro, Lightman & Eardley, 1976), and later for a class that became known as ADAFs (Advection Dominated Accretion Flows, Ichimaru, 1977; Narayan & Yi, 1994, 1995; Abramowicz et al., 1995). In this case, the cooling is negligible, either because the optical depth is large, or because the density is so low that the radiative efficiency is extremely small. The dissipated energy is advected with the flow and may be absorbed by the central object. Such flows are geometrically thick, with scale heights $`Hr`$. What little radiation does emerge from them arises, as in the thin disks, through various electromagnetic processes. In either of the above scenarios, accretion (and the accompanying radiation) is usually thought to be limited by the Eddington rate, a self–regulatory balance imposed by Newtonian gravity and radiation pressure. The standard argument gives a maximum luminosity $`L_{\mathrm{Edd}}=1.3\times 10^{38}(M/M_{\mathrm{}})`$ erg s<sup>-1</sup>. Although this may not be strictly the case in reality — as in the current argument concerning the nature of ULXs observationally (Rappaport, Podsiadlowski & Pfahl, 2005; King & Dehnen, 2005), and also quite general theoretical considerations (Abramowicz, 2004) — it does exhibit the qualitative nature of the effect of radiation pressure on accreting plasma, in the limit of large optical depth. The problem is circumvented (or at least deferred by nearly sixteen orders of magnitude in luminosity) if the main cooling agent is emission of neutrinos, instead of photons. This regime requires correspondingly large accretion rates, of the order of one solar mass per second, and is termed hypercritical accretion (Chevalier, 1989). It is important, for example, in the context of post–supernova fallback accretion onto a proto–neutron star. In such a situation, the densities and temperatures are so large ($`\rho 10^{12}`$ g cm<sup>-3</sup>, $`T10^{11}`$ K) that photons are completely trapped, and energetic neutrinos are emitted in large amounts, cooling the gas and allowing accretion to proceed. The reader may wish to consult the very clear introduction in this context given by Houck & Chevalier (1991). When the stellar envelope experiences fallback onto the central object, the system may be considered to be nearly in a steady state, because of the timescales involved. There are other, more violent situations, relevant to the study of gamma-ray bursts (GRBs) in which the assumption of steady state is not justified because of a time–varying mass and energy supply, and which therefore require a dynamical analysis for their proper description. In general, whether the system cools via electromagnetic radiation or neutrinos, analytic steady state solutions are found by assuming that the cooling is either efficient (thin disk) or inefficient (thick disk). Over the last several years, many groups have considered these cases in the neutrino cooled regime (Popham, Woosley & Fryer, 1999; Narayan, Piran & Kumar, 2001; Kohri & Mineshige, 2002; DiMatteo, Perna & Narayan, 2002; Yokosawa et al., 2004) and established the general features of the solutions. The real solution may be very different, in particular (and in what concerns us for the following) because of the previous history of the gas that constituted the accretion disk and how the disk formed. It turns out this may be quite important in the case of GRBs. Numerical analysis of the evolution of the disk with no assumptions concerning the steady state has been carried out by Setiawan et al. (2004). They have reported a three dimensional calculation with detailed microphysics. Unfortunately, computational limitations do not allow one to continue such calculations for more than approximately 50 ms, whereas the typical short GRB lasts 0.2 s. The sources of GRBs are now established (at least based on observations of X–ray, optical and radio counterparts) to lie at cosmological distances, with redshifts $`z14`$ (van Paradijs, Kouveliotou & Wijers, 2000, and references therein). At such scales, the absolute energetics of each event is approximately $`10^{52}`$ erg, assuming isotropic emission. Apparent collimation inferred from achromatic breaks in the afterglow light curve may reduce this to $`10^{50}`$ erg, depending on the particular event (Frail et al., 2001; Panaitescu & Kumar, 2001; Berger, Kulkarni & Frail, 2003). Although one generally considers two classes of GRBs based on duration (shorter and longer than about 2 s, Kouveliotou et al., 1993), all afterglows to date (and the corresponding inferences) come from long bursts (although see Lazzati, Ramirez-Ruiz & Ghisellini, 2001). Strong evidence in favor of a SN/GRB association (Kulkarni et al., 1998; Stanek et al., 2003; Hjorth et al., 2003) now comes from several events (albeit all relatively nearby), hinting at an underlying physical connection between the two phenomena. This is the central basis of the collapsar model (Woosley, 1993; MacFadyen & Woosley, 1999), in which rotation of the pre–supernova star allows for the creation of a massive accretion disk, fed by the infalling stellar envelope. The GRB is then powered for a fallback time, which can be long enough to account for the observed durations. In the case of short bursts there is less evidence to go on, but one possibility is that they arise from compact binary mergers, with various combinations of black holes and neutron stars in the binary (Lattimer & Schramm, 1976; Paczyński, 1986; Eichler et al., 1989; Paczyński, 1991; Narayan, Paczyński & Piran, 1992) such as in the first binary pulsar to be discovered (Hulse & Taylor, 1975). The tidal disruption of one star by the other gives rise to an accretion disk which powers the GRB. A black hole is either present from the start (as a variation of the Hulse–Taylor system) or is produced by the collapse of a supramassive neutron star shortly after the merger itself. In this case there is no external agent feeding the accretion disk, and thus the event is over roughly on an accretion timescale (which would be on the order of one second). Clearly, investigating either of these scenarios requires time–dependent, multidimensional calculations with detailed microphysics (Ruffert, Janka & Schäfer, 1996; Kluźniak & Lee, 1998; Ruffert & Janka, 1999; Lee, 2001; Rosswog & Ramirez–Ruiz, 2002; Rosswog, Ramirez–Ruiz & Davies, 2003; Rosswog, Speith, & Wynn, 2004). The energy released by accretion is then transformed into a relativistic outflow which produces the gamma ray burst (see Mészáros, 2002; Zhang & Mészáros, 2004; Piran, 2004, for reviews). We note that either scenario would produce a distinct gravitational wave signal, which could in principle be detected by interferometric systems such as LIGO in the future, thus establishing the nature of the progenitor system without a doubt. The regime in which the gas lies is unlike any other commonly encountered in astrophysics. The high densities and temperatures lead to photodisintegration of the nuclei and the establishment of nuclear statistical equilibrium (NSE). Furthermore, neutronization becomes important in the innermost regions of the flow and weak interactions determine the composition, with the electron fraction falling substantially below 1/2, and the gas correspondingly becoming neutron–rich. If this composition is somehow frozen and transported out of the gas and into an outflow, interesting nucleosynthesis of heavy elements could occur (Qian & Woosley, 1996; Pruet, Woosley & Hoffman, 2003; Pruet, Thompson & Hoffman, 2004). Realistic physics input of this kind allow us to obtain more reliable estimates of the actual energy released from the disk, and potentially available to power a GRB. An added complication, which affects the composition (Beloborodov, 2003) is that even if the main cooling mechanism is neutrino emission, these are not entirely free to leave the system, as scattering (mainly off free nucleons) is important enough to suppress the emission in the dense inner disk. We find that in fact the opaqueness of the material may lead to convection through the establishment of a composition gradient that does not satisfy the classical requirements for stability. To our knowledge, this is the first time that this has been addressed in this context. In previous work we initially studied the merger process for black hole–neutron star binaries in three dimensions, paying particular attention to the structure of the accretion disks that formed as a result of the tidal disruption (Lee, 2001). Follow–up work used the results of these simulations, mapped to two dimensions in azimuthal symmetry, as initial conditions with simple (ideal gas) input physics, and no realistic cooling included (Lee & Ramirez–Ruiz, 2002). A first approximation at a realistic equation of state and the effects of neutrino opacities was reported more recently (Lee, Ramirez–Ruiz & Page, 2004). In this paper we improve upon our earlier results in three important ways, that make them more realistic, and particularly relevant to the study of GRBs. First, we use a much more detailed equation of state, appropriate for the actual physical conditions found in post–merger accretion disks. This includes a relativistic electron gas of arbitrary degeneracy, radiation pressure and an ideal gas of free nucleons and $`\alpha `$ particles. The composition is determined self–consistently by considering weak interactions in two different regimes: neutrino–opaque and neutrino–transparent. Second, we include neutrino emission as the main source of cooling, by considering the relevant reaction rates (electron and positron capture onto free nucleons, bremsstrahlung, pair annihilation and plasmon decays), taken from tables and fitting formulae valid over wide ranges of temperature and density whenever possible. Third, we compute an approximate optical depth for the fluid to neutrinos, using scattering off free nucleons and $`\alpha `$ particles as a source of opacity. The emission rates and pressure due to neutrinos are then suppressed and enhanced, respectively, by an appropriate factor. We begin this paper by a presentation of the physical conditions likely to occur in post–merger accretion disks in §2, immediately following the coalescence of the two compact objects. Our results follow in §3, and a discussion of these, and the implications they have for GRBs is presented in §4. ## 2 Post–merger accretion disks The accurate study of the dynamical evolution of the accretion disks requires knowledge of the physical conditions within them, and the use of an appropriate equation of state. Below we give a general overview of the physical conditions in the disk, followed by a detailed presentation of the equation of state and the effects of neutrinos. ### 2.1 Physical conditions The accretion structures that are formed as a result of the merger of two neutron stars or the tidal disruption of a neutron star by the black hole become azimuthally symmetric fairly quickly, within a few tens of milliseconds (Ruffert, Janka & Schäfer, 1996; Lee, 2001). These disks typically contain a few tenths of a solar mass, and are small, with the bulk of the mass being contained within 200 km of the black hole (which harbors about 3–5 solar masses). The disks are dense ($`10^9\rho \text{[g cm}\text{-3}\text{]}10^{12}`$), with high internal energies ($`10^{10}T\text{[K]}10^{11}`$) (see e.g., Ruffert & Janka, 1999; Rosswog et al., 1999; Lee & Ramirez–Ruiz, 2002). In fact, the temperature is high enough that nuclei become photodisintegrated and there is a mixture of $`\alpha `$ particles, free neutrons and protons, electrons and positrons. The timescale for $`\beta `$–equilibrium is given by $`t_\beta (\sigma _{Ne}n_\mathrm{N}c)^1`$, where $`\sigma _{Ne}`$ is the cross section for $`e^\pm `$ capture onto free nucleons and $`n_\mathrm{N}`$ is the number density of free nucleons (see e.g., Shapiro & Teukolsky, 1983). Under these conditions, $`t_\beta 2\times 10^4`$ s, which is much shorter than the accretion timescale $`t_{acc}`$ of a fluid element, so that weak interactions determine the composition. Photons are trapped and are advected with the flow. For neutrinos the situation is more complicated, since the optical depth is of order unity. In the outer regions of the disk, they escape freely, whereas for densities larger than $`10^{11}`$ g cm<sup>-3</sup> they undergo diffusion on a timescale $`t_{\mathrm{dif},\nu }30`$ ms. ### 2.2 The equation of state and the composition of the fluid. For simplicity of presentation, we will first assume that all nucleons are free (no $`\alpha `$ particles are present), and include the necessary corrections subsequently. Under the conditions described above, the temperature is high enough for electron–positron pair creation. The number of pairs thus produced is very sensitive to the degeneracy of the electrons. In fact the number density of pairs is suppressed exponentially with increasing degeneracy, because of Fermi blocking. In previous work (Popham, Woosley & Fryer, 1999; Narayan, Piran & Kumar, 2001; Kohri & Mineshige, 2002; DiMatteo, Perna & Narayan, 2002; Lee, Ramirez–Ruiz & Page, 2004) the assumption of full degeneracy has been made for simplicity (note also that in some cases the presence of electron–positron pairs has been assumed while at the same time retaining a degeneracy pressure term, which is inconsistent). Here we take a different approach, using an exact expression for the pressure as a function of the temperature and the chemical potential, valid for arbitrary degeneracy in the limit of relativistic electrons (i.e., $`\rho 10^6`$ g cm<sup>-3</sup>), due to Blinnikov, Dunina–Barkovskaya & Nadyozhin (1996), namely: $$P_e=\frac{1}{12\pi ^2(\mathrm{}c)^3}\left[\eta _e^4+2\pi ^2\eta _e^2(kT)^2+\frac{7}{15}\pi ^4(kT)^4\right].$$ (1) The number densities of electrons and positrons are related by: $$\frac{\rho Y_e}{m_u}=n_{}n_+=\frac{1}{3\pi ^2(\mathrm{}c)^3}[\eta _e^3+\eta _e\pi ^2(kT)^2]$$ (2) where $`Y_e`$ is the electron fraction. The chemical potential of species $`i`$ is denoted by $`\eta _i`$ throughout. This expression reduces to the well–known limits when the temperature is low ($`kT\eta _e`$, which gives $`P\rho ^{4/3}`$, appropriate for a cold relativistic Fermi gas) and when it is high ($`kT\eta _e`$, which gives $`PT^4`$, when the pressure comes from relativistic electron–positron pairs). The full equation of state then reads: $$P=P_{\mathrm{rad}}+P_{\mathrm{gas}}+P_e+P_\nu ,$$ (3) where $$P_{\mathrm{rad}}=\frac{aT^4}{3},$$ (4) $$P_{\mathrm{gas}}=\frac{\rho kT}{m_u},$$ (5) and $`P_\nu `$ is the pressure due to neutrinos (discussed below). Here $`a`$ is the radiation constant, $`k`$ is Boltzmann’s constant and $`m_u=1.667\times 10^{24}`$ g is the atomic mass unit. Since the presence of pairs is automatically taken into account in the expression for $`P_e`$, there is no alteration to the numerical factor $`1/3`$ in the expression for $`P_{\mathrm{rad}}`$. For the conditions encountered in the accretion disks presented below, gas pressure dominates at the 80% level over the other terms. Regarding the photons, since the temperature is $`T5`$ MeV, the peak in the blackbody spectrum is at $`14`$ MeV. On the other hand, the plasma frequency, $`\omega _p`$, corresponds to $`T_p0.5`$ MeV. Thus the standard expression for radiation pressure may be considered accurate, and plasma effects negligible as far as the photons are concerned. The computation of the electron fraction and the chemical potential of electrons follows from the assumption of $`\beta `$–equilibrium between neutrons, protons and electrons, and the condition of charge neutrality. As noted by (Beloborodov, 2003), a distinction needs to be made to determine the equilibrium composition depending on the optical depth of the material (the determination of the opacities is presented in § 2.4). If it is transparent to its own neutrino emission, we fix equilibrium by equating the capture rates of electrons and positrons onto protons and neutrons respectively. For mild degeneracy, as is the case here, this leads to the following expression for the electron fraction as a function of the temperature and the electron chemical potential (Beloborodov, 2003): $$Y_e=\frac{1}{2}+0.487\left(\frac{Q/2\eta _e}{kT}\right),$$ (6) where $`Q=(m_nm_p)c^21.29`$ MeV. If, however, the material is opaque, and the neutrinos are allowed to diffuse out on a timescale shorter than the accretion timescale, then we may write $$\eta _e+\eta _p=\eta _n,$$ (7) for the equilibrium composition<sup>1</sup><sup>1</sup>1This expression neglects the neutrino chemical potential, $`\eta _\nu `$, and is strictly valid only when there are equal numbers of neutrinos and anti–neutrinos. Otherwise the electron pressure, bulk viscosity and Joule–like heating due to the induced lepton current will be affected (Burrows, Mazurek & Lattimer, 1981; Socrates et al., 2004). We neglect all these corrections in the present treatment.. Since the nucleons are not degenerate, we may use Maxwell–Boltzmann statistics to describe their distribution function, and obtain $$\frac{n_p}{n_n}=\mathrm{exp}[(Q\eta _e)/kT]$$ (8) for the ratio of proton to neutron number densities. Further, with $`Y_e=n_p/(n_p+n_n)`$ we arrive at $$\frac{1Y_e}{Y_e}=\mathrm{exp}[(\eta _eQ)/kT].$$ (9) We now have the three equations (2), (3) and (6) or (9) for the three functions $`T,Y_e,\eta _e`$ and so the system is closed. To allow for a transition from the optically thin to optically thick regime, in fact we solve these in a combined form, weighted by factors $`f(\tau _\nu )=\mathrm{exp}[\tau _\nu ]`$ or $`g(\tau _\nu )=(1\mathrm{exp}[\tau _\nu ])`$. As a practical matter, the internal energy per unit mass, $`u`$, is used instead of the pressure to solve equation (3), since its variation in time is what is determined in the code, using the First Law of Thermodynamics. We finally include the effects of incomplete photodisintegration of $`\alpha `$ particles by using nuclear statistical equilibrium for the three species $`(n,p,\alpha )`$ to fix the mass fraction of free nucleons as (Qian & Woosley, 1996): $$X_{\mathrm{nuc}}=22.4\left(\frac{T}{[10^{10}\text{K}]}\right)^{9/8}\left(\frac{\rho }{10^{10}\text{g cm}\text{-3}}\right)^{3/4}\mathrm{exp}(8.2[10^{10}\text{K}]/T).$$ (10) Whenever this expression results in $`X_{\mathrm{nuc}}>1`$ we set $`X_{\mathrm{nuc}}=1`$. The corresponding alterations in the previous derivation are simple, and lead to the full set of equations which we write for each gas element and at each time. Note that these do not allow for an explicit solution, so an iterative scheme is used in each instance. The internal energy per unit mass is $$u=3\frac{P_e+P_{\mathrm{rad}}+P_\nu }{\rho }+\frac{1+3X_{\mathrm{nuc}}}{4}\frac{3kT}{2m_u},$$ (11) $`\beta `$–equilibrium gives $$Y_e=(1X_{\mathrm{nuc}})/2+X_{\mathrm{nuc}}([\frac{1}{2}+0.487\left\{\frac{Q/2\eta _e}{kT}\right\}]f(\tau _\nu )+[1+\mathrm{exp}(\{\eta _eQ\})/kT)]g(\tau _\nu )),$$ (12) and charge neutrality implies $$\frac{\rho Y_e}{m_u}=\frac{1}{3\pi ^2(\mathrm{}c)^3}[\eta _e^3+\eta _e\pi ^2(kT)^2],$$ (13) as given by equation (2). The modifications to the equation for $`\beta `$ equilibrium simply reflect the fact that if the fluid is composed primarily of $`\alpha `$ particles, there will be one neutron per proton, and hence $`Y_e1/2`$. ### 2.3 Neutrino emission and photodisintegration losses Several processes contribute to the emission of neutrinos. First, we take into account electron and positron capture onto free nucleons using the tables of Langanke & Martínez–Pinedo (2001). This is an improvement over calculations done previously (Popham, Woosley & Fryer, 1999; Narayan, Piran & Kumar, 2001; Kohri & Mineshige, 2002; DiMatteo, Perna & Narayan, 2002; Lee, Ramirez–Ruiz & Page, 2004), where assumptions concerning the degeneracy were made. The rates are thus accurate over the entire disk, whether the degeneracy is significant or not. The tables cover the following ranges: $`1\mathrm{log}(\rho Y_e\text{[g cm}\text{-3}\text{]})11`$; $`7\mathrm{log}T\text{[K]}11`$. A bilinear interpolation in the $`\mathrm{log}\rho Y_e\mathrm{log}T`$ plane is performed using the table to obtain the cooling rate for a given mass element. The result is then multiplied by the mass fraction of free nucleons, $`X_{\mathrm{nuc}}`$, since we are not considering capture of electrons and protons by $`\alpha `$ particles. Second, the annihilation of electron–positron pairs is a source of thermal neutrinos, and the corresponding cooling rate is computed using the fitting functions of Itoh et al. (1996). These cover the range $`9\mathrm{log}(\rho \text{[g cm}\text{-3}\text{]})12`$ in density and $`10\mathrm{log}T\text{[K]}11`$ in temperature. Finally, nucleon–nucleon bremsstrahlung, $`\dot{q}_{ff}`$, and plasmon decay, $`\dot{q}_{\mathrm{plasmon}}`$, are considered, with rates given by: $$\dot{q}_{ff}=1.5\times 10^{33}T_{11}^{5.5}\rho _{13}^2\text{erg s}\text{-1}\text{ cm}\text{-3},$$ (14) (Hannestad & Raffelt, 1998) and $$\dot{q}_{\mathrm{plasmon}}=1.5\times 10^{32}T_{11}^9\gamma _p^6\mathrm{exp}(\gamma _p)(1+\gamma _p)\left(2+\frac{\gamma _p^2}{1+\gamma _p}\right)\text{erg s}\text{-1}\text{ cm}\text{-3},$$ (15) where $`\gamma _p=5.5\times 10^2\sqrt{(\pi ^2+3[\eta _e/kT]^2)/3}`$ (Ruffert, Janka & Schäfer, 1996). For the conditions encountered in the disk, capture by nucleons completely dominates all other processes, and is the main source of cooling. Finally, the creation and disintegration of $`\alpha `$ particles leads to a cooling term in the energy equation given by: $$\dot{q}_{\mathrm{phot}}=6.8\times 10^{18}\frac{dX_{\mathrm{nuc}}}{dt}\text{erg s}\text{-1}\text{ cm}\text{-3}.$$ (16) ### 2.4 Neutrino opacities The material in the disk is dense enough that photons are completely trapped, and advected with the flow. For neutrinos however, the situation is more complicated, since the outer regions are transparent to them, while the inner portions are opaque. The main source of opacity is scattering off free nucleons and $`\alpha `$ particles, with cross–sections given by (Tubbs & Schramm, 1975; Shapiro & Teukolsky, 1983) $$\sigma _\mathrm{N}=\frac{1}{4}\sigma _0\left(\frac{E_\nu }{m_ec^2}\right)^2,$$ (17) and $$\sigma _\alpha =\sigma _0\left(\frac{E_\nu }{m_ec^2}\right)^2\left[(4\mathrm{sin}^2\theta _w)\right]^2,$$ (18) where $`\sigma _0=1.76\times 10^{44}`$ cm<sup>2</sup> and $`\theta _w`$ is the Weinberg angle. Since capture processes dominate the neutrino luminosity, we assume that the mean energy of the neutrinos is roughly equal to the Fermi energy of the partially degenerate electrons $`E_\nu \frac{5}{6}\eta _e=43(Y_e\rho _{12})^{1/3}`$ MeV, where $`\rho _{12}=\rho /10^{12}`$g cm<sup>-3</sup>. So the mean free path is $`l_\nu =1/(n_\mathrm{N}\sigma _\mathrm{N}+n_\alpha \sigma _\alpha )`$, $`n_\mathrm{N}`$ and $`n_\alpha `$ being the number densities of free nucleons and $`\alpha `$ particles respectively. To compute an optical depth, we take $`\tau _\nu =H/l_\nu `$, with $`H`$ a typical scale height of the disk. Inspection of the disk shape (through the density contours in the inner regions) and an estimation of the scale height with $`H\rho /\rho `$ suggests that $`H\kappa r`$, with $`\kappa `$ a constant of order unity, so that our final, simplified expression for the optical depth reads $`\tau _\nu =\kappa r/l_\nu `$. With the above expressions for the cross section as a function of density and electron fraction, and defining $`r_7=r/10^7\text{cm}`$ we arrive at $$\tau _\nu =186.5\kappa \rho _{12}^{5/3}Y_e^{2/3}r_7\left[X_{\mathrm{nuc}}+3.31(1X_{\mathrm{nuc}})/4\right].$$ (19) The term which depends on $`X_{\mathrm{nuc}}`$ reflects the two–species composition we have assumed. Its influence is limited, since when $`X_{\mathrm{nuc}}`$ varies from 0 to 1, the optical depth is altered by a factor 1.2, all other values being equal. In practice, and for the compositions found in the disks, this expression corresponds to having an optical depth of unity at approximately $`10^{11}`$g cm<sup>-3</sup>, which can be thought of as a neutrino–surface, where the last scattering occurs and neutrinos leave the system. Accordingly, we modify the expressions for the neutrino emission rates, suppressing it in the opaque regions, and enhancing the pressure through $$\left(\frac{du}{dt}\right)_\nu =\left(\frac{du}{dt}\right)_0\mathrm{exp}(\tau _\nu ),$$ (20) and $$P_\nu =\frac{7}{8}aT^4[1\mathrm{exp}(\tau _\nu )],$$ (21) where $`(du/dt)_0=\dot{q}_i/\rho `$ is the unmodified energy loss rate (in erg g<sup>-1</sup> s<sup>-1</sup>) calculated from the rates given in § 2.3. The total neutrino luminosity (in erg s<sup>-1</sup>) is then computed according to $$L_\nu =\rho ^1\left[\dot{q}_{ff}+\dot{q}_{\mathrm{plasmon}}+\dot{q}_{\mathrm{pair}}+\dot{q}_{\mathrm{cap}}\right]\mathrm{exp}(\tau _\nu )𝑑m.$$ (22) ### 2.5 Numerical method and initial conditions. For the actual hydrodynamical evolution calculations, we use the same code as in previous work (Lee & Ramirez–Ruiz, 2002), and refer the reader to that paper for the details. This is a two dimensional (cylindrical coordinates $`(r,z)`$ in azimuthal symmetry) Smooth Particle Hydrodynamics code (Monaghan, 1992), modified from our own 3D version used for compact binary mergers (Lee, 2001). The accretion disk sits in the potential well of a black hole of mass $`M_{\mathrm{BH}}`$, and which produces a Newtonian potential $`\mathrm{\Phi }=GM_{\mathrm{BH}}/r`$. Accretion is modeled by an absorbing boundary at the Schwarzschild radius $`r_{Sch}=2GM_{\mathrm{BH}}/c^2`$, and the mass of the hole is updated continuously. The transport of angular momentum is modeled with an $`\alpha `$ viscosity prescription, including all components of the viscous stress tensor (not only $`t_{r\varphi }`$). The self–gravity of the disk is neglected. The main modifications from the previous version are in the implementation of a new equation of state §2.2, neutrino emission, §2.3 and the treatment of neutrino optical depths, §2.4. As done previously (Lee & Ramirez–Ruiz, 2002; Lee, Ramirez–Ruiz & Page, 2004), the initial conditions are taken from the final configuration of 3D calculations of black hole–neutron star mergers, after the accretion disk that forms through tidal disruption of the neutron star has become fairly symmetric with respect to the azimuthal coordinate. We show in Table 1 the parameters used in each of the dynamical runs included in this paper. ## 3 Results As the dynamical evolution calculations begins, the disks experience an initial transient, which is essentially numerical in origin, and is due to the fact that the configuration is not in strict hydrostatic equilibrium (recall that it is obtained from azimuthally averaging the results of 3D calculations). This transient, and its effects, are negligible, and it is essentially over in a hydrodynamical timescale ($`2`$ ms). After that, the disk proceeds to evolve on a much longer timescale, determined by accretion onto the central black hole. We first present the details of the spatial structure of the disk due to fundamental physical effects, and then proceed to show the temporal evolution and associated transitions. ### 3.1 Disk structure The fundamental variable affecting the instantaneous spatial structure of the disk is the optical depth to neutrinos, $`\tau _\nu `$, since it determines whether the fluid cools efficiently or not. All quantities show a radical change in behavior as the threshold $`\tau _\nu =1`$ is passed. Figure 1 shows color–coded contours in a meridional slice of the thermodynamical variables, while Figure 2 displays the run of density and entropy per baryon along the equator, $`z=0`$, where the density and temperature are highest (we will refer here mainly to run a2M, unless noted otherwise). The density increases as one approaches the black hole, varying as $`\rho r^5`$ in the outer regions, where $`\tau _\nu 1`$. The critical density for opaqueness is reached at $`r_{}10^7`$ cm, and for $`r<r_{}`$, $`\rho r^1`$. Since the energy that would otherwise be lost via neutrino emission remains in the disk if the cooling is suppressed by scattering, the density does not rise as fast. In fact, one can note this change also by inspecting the scale height $`HP/P`$, which scales as $`H/rr`$ for $`r>r_{}`$ and $`H/r`$ const for $`r<r_{}`$ (and thus in terms of surface density the change is from $`\mathrm{\Sigma }r^4`$ to $`\mathrm{\Sigma }`$ const at $`r_{}`$). The opaque region of the disk remains inflated to a certain extent, trapping the internal energy it holds and releasing it only on a diffusion timescale. At densities $`\rho 10^{12}`$ g cm<sup>-3</sup> the optical depth can reach $`\tau _\nu 100`$, and thus the suppression of cooling is dramatic, as can be seen in Figure 1, where the contours of $`\tau _\nu `$ and $`\dot{q}`$ are shown. The flattening of the entropy profile is related to the change in composition in the optically thick region and the occurrence of convection (see § 3.2 below). As the density rises in the inner regions of the disk, the electron fraction $`Y_e`$ initially drops as neutronization becomes more important, with the equilibrium composition being determined by the equality of electron and positron captures onto free neutrons and protons (see equation (6) and the discussion preceding it). The lowest value is reached at $`rr_{}`$, where $`Y_e0.03`$. Thereafter it rises again, reaching $`Y_e0.1`$ close to the horizon. Thus flows that are optically thin everywhere will reach a higher degree of neutronization close to the black hole than those which experience a transition to the opaque regime. The numerical values for the electron fraction at the transition radius and at the inner boundary are largely insensitive to $`\alpha `$, as long as the transition does occur. The baryons in the disk are essentially in the form of free neutrons and protons, except at very large radii and low densities ($`r>5\times 10^7`$ cm and $`\rho <5\times 10^6`$g cm<sup>-3</sup>) where $`\alpha `$ particles form. Figure 3 shows the region in the density–temperature plane where the fluid lies, for runs a2M and a1M at two different times (and color coded according to the electron fraction). Most of the gas lies close to the line determining the formation of Helium nuclei (see equation 10), but does not cross over to lower temperatures. The reason for this is that the energy that would be released by the creation of one Helium nucleus (28.3 MeV) would not leave the disk (recall that we consider neutrino emission as the only source of cooling) and thus immediately lead to the photodisintegration of another $`\alpha `$ particle. An equilibrium is thus maintained in which the gas is close to Helium synthesis, but this does not occur. In the opaque regime, the gas moves substantially farther from the transition line to Helium, since it cannot cool efficiently and, as mentioned before, the density does not rise as quickly. Figure 3 also shows clearly why one must make use of an equation of state which takes into account properly the effects of arbitrary degeneracy of the electrons and positrons. The solid straight line marks the degeneracy temperature as a function of density, given by $`kT=7.7\rho _{11}^{1/3}`$ MeV. The disk straddles this line, with a degeneracy parameter $`\eta _e/kT24`$ in the inner regions, and $`\eta _e/kT1`$ at lower densities. Thus making the approximation that the electrons in the flow are fully degenerate is not accurate. Modeling the accretion flow in two dimensions $`(r,z)`$, without the assumption of equatorial symmetry allows one to solve clearly for the vertical motions in the disk, something which is not possible when considering vertically–integrated flows. It was pointed out by Urpin (1984) that a standard $`\alpha `$ viscosity could lead to meridional flows in which $`v_r`$ changed sign as a function of height above the mid plane, leading to inflows as well as outflows. Kita (1995), Kluźniak & Kita (2000) and Regev & Gitelman (2002) later considered a similar situation, and found that for a range of values in $`\alpha `$, the gas flowed inward along the surface of the disk, and outward in the equatorial regions. In previous work (Lee & Ramirez–Ruiz, 2002), we found this solution in disk flow, with large scale circulations exhibiting inflows and outflows for high viscosities ($`\alpha 0.1`$), and small–scale eddies at lower values ($`\alpha 0.01`$). How vertical motions affect the stability of accretion disks and may lead to the transport of angular momentum is a question that has become of relevance in this context (Arlt & Urpin, 2004). Here we show in Figure 4 the velocity field and magnitude of meridional velocity for run a2M. The small–scale eddies are clearly visible, with their strength usually diminishing as the disk is drained of matter and the density drops. The effect of a different value of $`\alpha `$ can be seen in the top row of Figure 5, where the comparison between $`\alpha =0.1`$ and $`\alpha =0.01`$ is made. The large, coherent lines of flow aimed directly at the origin for $`\alpha =0.1`$ correspond to inflow, while the lighter shades along the equator at larger radii show outflowing gas. The instantaneous structure of the disk concerning the density, temperature and composition profiles is largely independent of the viscosity, as long as there is enough mass to produce the optically thin/optically thick transition. Even in the optically thin regime, the structure cannot be determined analytically following the standard arguments applied to thin, cool disks of the Shakura–Sunyaev type. The reason for this is the following: a central assumption in the standard solution is that the disk is cool, i.e., $`kT/m_pGM_{\mathrm{BH}}/r`$. This comes from the requirement that all the energy dissipated by viscosity, $`\dot{q}_\alpha `$ be radiated away efficiently, and produce the observed flux, $`F`$. In our case, the disk is most certainly not cool, as it originates from the tidal disruption of a neutron star by a black hole (or possibly the coalescence of two neutron stars, and subsequent collapse of the central mass to a black hole). The gas that constitutes the disk is dynamically hot because it was in hydrostatic equilibrium in a self–gravitating configuration, where $`UW`$ and has not been able to release this internal energy (the merger process itself may lead to additional heating). A further deviation from the standard solution is that the pressure support in the disk leads to a rotation curve that is slightly sub–Keplerian. The dissipation by viscosity is in fact smaller than the cooling rate over much of the disk, and the released luminosity comes from a combination of viscous dissipation and the store of internal energy given to the disk at its conception. ### 3.2 Lepton–driven convection The classical requirement for convective instability in the presence of entropy and composition gradients, as well as rotation, is the Solberg–Hoiland criterion (Tassoul, 1978): $$N^2+\omega _r^2<0,$$ (23) where $`\omega _r^2=4\mathrm{\Omega }^2+rd\mathrm{\Omega }^2/dr`$ is the radial epicyclic frequency ($`\mathrm{\Omega }`$ being the angular velocity), and $$N^2=\frac{g}{\gamma }\left[\frac{1}{P}\left(\frac{P}{s}\right)_{Y_e}\frac{ds}{dr}+\frac{1}{P}\left(\frac{P}{Y_e}\right)_s\frac{dY_e}{dr}\right]$$ (24) is the Brunt–Väisäla frequency (see Lattimer & Mazurek, 1981; Thompson, Quataert & Burrows, 2005). The adiabatic index $`\gamma =d\mathrm{ln}P/d\mathrm{ln}\rho 5/3`$ in our case, since $`P_{gas}`$ is the most important contribution to the total pressure. Thus a region may be convectively unstable because of a composition gradient, or an entropy gradient, or both. Strictly speaking, here one should consider the total lepton fraction $`Y_l`$. We do not consider the transport of neutrinos in detail, as already mentioned (see § 2.4), and do not calculate explicitly the contribution of neutrinos to this $`Y_l`$. For what follows, we will thus simply take $`Y_e`$ to represent the behavior of the full $`Y_l`$. The origin of convection in the lepton inversion zone can be understood as follows (Epstein, 1979). Consider a fluid element in the lepton inversion zone that is displaced in the outward direction and then comes to pressure equilibrium with its surroundings. The displaced element, which is lepton rich relative to its new surroundings, attains the ambient pressure at a lower density than the surrounding fluid, because the pressure depends directly on the lepton number, and thus tends to drift outwards. By the same token, an inwardly displaced fluid element in the lepton inversion zone is depleted in leptons relative to its new surroundings and thus tends to sink <sup>2</sup><sup>2</sup>2This is similar to the thermosolutal convection which would occur in a normal star if there were an inwardly decreasing molecular weight gradient or composition inversion (Spiegel, 1972).. When temperature gradients are allowed for, the displaced fluids, in general, have temperatures which differ from those of their surroundings. These variations tend to promote stability or instability depending on whether the existing temperature gradient is less than or greater than the adiabatic gradient, respectively. In the present scenario then, the transition to the optically thick regime leads to instability, because initially $`ds/dr<0`$, and also $`dY_e/dr<0`$ (see also Figure 2). The entropy profile is then flattened by efficient convective mixing, and the lepton gradient inversion is due to the different way in which the composition is determined through weak interactions once the neutrinos become trapped. In a dynamical situation such as the one treated here, convection tends to erase the gradients which give rise to it. Accordingly, the sum of the corresponding terms in equation (23) tends to zero once the simulation has progressed and convection has become established in the inner disk (note that in our case the rotation curve there is sub–Keplerian, with $`\mathrm{\Omega }r^{8/5}`$, so that $`\mathrm{\Omega }`$ and $`\omega _r`$ are not equal). This is analogous to what occurs in proto neutron stars following core collapse (Epstein, 1979; Burrows & Lattimer, 1986; Thompson & Duncan, 1993) and has actually been confirmed in numerical simulations of such systems (Janka & Müller, 1996), where the convection leads to strong mixing. For a neutrino–driven convective luminosity $`L_{\mathrm{con}}`$ at radius $`R`$ and density $`\rho `$, the convective velocity may be estimated by standard mixing length theory as (see e.g., Thompson & Duncan, 1993) $$v_{\mathrm{con}}3.3\times 10^8\left(\frac{L_{\mathrm{con}}}{10^{52}\mathrm{erg}\mathrm{s}^1}\right)^{1/3}\left(\frac{\rho }{10^{12}\mathrm{g}\mathrm{cm}^3}\right)^{1/3}\left(\frac{R}{20\mathrm{km}}\right)^{2/3}\mathrm{cm}\mathrm{s}^1,$$ (25) where we have used the fact that gas pressure dominates in the fluid. The assumption of spherical symmetry implicit in this expression is not strictly met in our case, but it may provide us nevertheless with a useful guide. The overturn time of a convective cell is then given by $`t_{\mathrm{con}}l_p/v_{\mathrm{con}}10`$ ms, where $`l_p20`$ km is the mixing length, usually set to a pressure scale height. In our calculations, we see that the magnitude of the meridional velocity $`\|v\|=\sqrt{v_r^2+v_z^2}`$ decreases as $`r`$ decreases, then rises again as the opaque region is reached, reaching $`\|v\|10^8`$ cm s<sup>-1</sup>, at $`r20`$ km. The associated turnover times are thus $`t_{\mathrm{con}}l_p/\|v\|20`$ ms, in good agreement with the estimate made above. To isolate this effect from that due to viscosity (which generates the meridional circulations mentioned previously) we have performed a simulation in which $`\alpha =0`$ (run aIM in Table 1). This calculation shows essentially no accretion onto the black hole except for a small amount of gas transferred at early times (because the initial condition is not in strict equilibrium). The result concerning the profile of $`Y_e`$ as a function of radius is as we have described above, i.e., increasing neutronization as $`r`$ decreases, until the opaque region is reached, followed by an increase in $`Y_e`$ as the black hole boundary is approached. The revealing difference lies in the time evolution of this profile. As mentioned above, convection generates motions which tend to eliminate the composition gradient which drives it. In the presence of viscosity, matter is continuously transported radially, and the gradient is not erased entirely. Height integrated profiles of $`Y_e(r)`$ at 50 ms and 200 ms show a similar behavior at small radii. When viscosity is removed, the initial composition gradient is gradually softened until it disappears in the innermost regions (see Figure 7). The disk then has a nearly constant electron fraction along the equator for $`r<r_{}`$, with a sharp transition region leading to an increase with radius in the transparent regime. This may resemble the convection dominated accretion solution found analytically by Quataert & Gruzinov (2000) and numerically by several groups (Stone, Pringle & Begelman, 1999; Narayan, Igumenshchev & Abramowicz, 1999), where convection transports angular momentum inward, energy outward, and gives a radial profile in density $`r^{1/2}`$. We find a different power law, with $`\rho r^1`$. There are several factors which may account for this difference: the disk vs. spherical geometry; the initial condition with a large amount of internal energy; and the limited, but present cooling rate on the boundaries of the flow. ### 3.3 Neutrino diffusion effects Since the neutrinos are diffusing out of the fluid (the mean free path is small compared to the size of the system in the optically thick regime), one would expect a corresponding viscosity through $`\zeta _\nu ^{\mathrm{vis}}`$ (see, e.g. Burrows, Mazurek & Lattimer, 1981). This needs to be compared to the corresponding viscosity generated by the assumed $`\alpha `$ prescription for consistency. We may assume $`\zeta _\nu ^{\mathrm{vis}}=(1/3)U_\nu l_\nu /c`$, where $`U_\nu aT^4`$ is the energy density associated with neutrinos, and $`l_\nu `$ is their mean free path (see § 2.4). Thus $$\zeta _\nu ^{\mathrm{vis}}2\times 10^{20}\text{g cm}\text{-1}\text{ s}\text{-1},$$ (26) for conditions near the equatorial plane, $`z=0`$. This value will increase in the outer regions, since the mean free path becomes larger as the density drops. At the neutrino–surface, where $`\tau _\nu 1`$, the mean free path is $`l_\nu 10^7`$ cm, and so an averaged value of the viscosity over the entire neutrino–opaque region will result in an increase of about one order of magnitude over the estimate given in equation (26) (note however, that strictly speaking the diffusion approximation is no longer valid in the outer regions, so this must be interpreted with care). For the $`\alpha `$ prescription, $`\zeta _\alpha =\rho \alpha c_s^2/\mathrm{\Omega }`$. The rotation curve is not too far from Keplerian, and scaling this expression to typical values we find $$\zeta _\alpha =1.3\times 10^{22}c_{s,9}^2\alpha _2r_7^{3/2}M_4^{1/2}\rho _{10}\text{g cm}\text{-1}\text{ s}\text{-1}.$$ (27) The effects on angular momentum transport are thus a full two orders of magnitude below those arising from our viscosity prescription (for $`\alpha =10^2`$). To put it another way, a lower limit for the viscosity under these conditions (and in the optically thick portion of the disk) would be $`\alpha 10^4`$. An alternative analysis would be to consider the corresponding timescales induced by this viscosity. Since $`t_{\mathrm{vis}}R(\rho /\zeta )^2`$, smaller viscosities imply longer timescales, as expected. ### 3.4 Stability Aside from convection, several general criteria for stability can be analyzed in our case. Evidently, since we are performing dynamical calculations, any instability that arises will quickly lead to a change in structure. It is nevertheless instructive to consider the corresponding conditions within the disk. We first consider the Toomre criterion, comparing gravitational and internal energies, with $$Q_\mathrm{T}=\frac{\omega _rc_s}{\pi G\mathrm{\Sigma }},$$ (28) where $`\omega _r`$ is the local epicyclic frequency (essentially equal in this case to the angular frequency $`\mathrm{\Omega }`$), $`c_s`$ is the local sound speed and $`\mathrm{\Sigma }`$ is the surface density. We find $`Q_\mathrm{T}>1`$ in all cases throughout the calculations, and thus that the disks are stable in this respect (Figure 8a shows a typical profile of $`Q_\mathrm{T}[r]`$). Previous studies (Narayan, Piran & Kumar, 2001; Kohri & Mineshige, 2002) have shown that neutrino cooled disks are thermally unstable if radiation pressure dominates in the flow, and stable otherwise, although the effects of neutrino opacities were not considered. The criterion for stability can be written in this case as $$\left(\frac{d\mathrm{ln}Q^+}{d\mathrm{ln}T}\right)_r\left(\frac{d\mathrm{ln}Q^{}}{d\mathrm{ln}T}\right)_r,$$ (29) where $`Q^+`$ and $`Q^{}`$ are the volume heating and cooling rates. Essentially, in order to be stable the disk must be able to get rid of any excess internal energy generated by an increase in the heating rate. We show in Figure 8b typical values for $`Q^+`$ and $`Q^{}`$ as functions of the central (equatorial) temperature, for run a2M at $`t=50`$ ms. In the optically thin region the disks are thermally stable, because the pressure is dominated by the contribution from free nucleons and cooling is efficient. In the optically thick region the cooling is greatly suppressed and the criterion would indicate that the disk becomes thermally unstable. This is reasonable, since with optical depths $`\tau _\nu 10100`$, essentially no direct cooling takes place, and dissipation is not suppressed. This is why the disk is geometrically thick in the inner regions, as already described above (§ 3.1). The balance that gives a quasi–steady state structure is achieved mainly through the diffusion of neutrinos, since the diffusion timescale is shorter than the accretion timescale. ### 3.5 Disk evolution The evolution of the disk on long timescales is determined by the balance between two competing effects: on one hand, viscosity transports angular momentum outwards, matter accretes and the disk drains into the black hole on an accretion timescale $`t_{\mathrm{acc}}`$. The trend in this respect is towards lower densities and temperatures. On the other hand, cooling reduces pressure support and leads to vertical compression, increasing the density. The internal energy of the fluid is released on a cooling timescale, $`t_{\mathrm{cool}}`$. The presence of an optically thick region in the center of the disk limits the luminosity (dominated by electron and positron capture onto free nucleons) to $`\mathrm{few}\times 10^{53}`$erg s<sup>-1</sup>, and the initial internal energy is $`E_{\mathrm{int}}10^{52}`$ erg, so $`t_{\mathrm{cool}}0.1`$ s. If $`t_{\mathrm{acc}}>t_{\mathrm{cool}}`$, the disk will cool before its mass or internal energy reservoir is significantly affected by accretion onto the black hole. The maximum density (shown in Figure 9 for runs a1M, a2M and a3M) actually increases slightly due to vertical compression, and subsequently drops once mass loss through accretion dominates (the $`10`$ ms delay at the start, during which it is approximately constant, is simply the sound crossing time across the optically thick region of the disk). The accretion rate onto the black hole, $`\dot{M}_{\mathrm{BH}}`$, the accretion timescale $`t_{\mathrm{acc}}=\dot{M}_{\mathrm{BH}}/M_{\mathrm{disk}}`$ and the total neutrino luminosity $`L_\nu `$ are shown in Figures 10, 11 and 12. They all show the same qualitative behavior, remaining fairly constant (or changing slowly) for an accretion timescale, and abruptly switching thereafter. The accretion timescales are approximately 0.5 s and 5 s for $`\alpha =0.01,0.001`$ respectively. For high viscosity, $`\alpha =0.1`$, the transport of angular momentum is vigorous, and the black hole quickly accretes a substantial amount of mass ( $`0.16M_{\mathrm{}}`$ within the first 100 ms). The accretion timescale is $`t_{\mathrm{acc}}50`$ ms. The circulation pattern consists of large–scale eddies, with $`Hr`$. In fact there is essentially one large eddy on each side of the equatorial plane, with mass inflow along the surface of the disk, and an equatorial outflow. Part of the outflowing gas moves away from the equator and reverses direction close to the surface of the disk, contributing to the inflow. For an intermediate viscosity, $`\alpha =0.01`$, the intensity of the circulations is smaller, but also, the eddies are smaller, with several of them clearly occurring in the disk at once. The transport of angular momentum being less vigorous, the accretion rate is substantially smaller than in the previous case. Finally, for a yet lower viscosity, $`\alpha =0.001`$, the trend continues, and the eddies become smaller still. In these last two cases, angular momentum transport is so low that very high densities are maintained in the central regions of the disk for a large number of dynamical times, contrary to what is seen in the high viscosity case. A simple way to quantify how much of the mass flow is actually making it to the central black hole is to measure the fraction of mass at any given radius that is moving inwards, $`\|\dot{M}_{\mathrm{in}}\|/(\|\dot{M}_{\mathrm{in}}\|+\|\dot{M}_{\mathrm{out}}\|)`$, where we have divided the mass flow rate into two height–integrated parts, one with $`v_r<0`$ and another with $`v_r>0`$. For example, for run a2M at $`t=100`$ ms and $`r=r_{}`$ it is approximately 1/3. The ratio tends to unity only in the innermost regions of the disk, and shows the true black hole mass accretion rate, plotted in Figure 10 for the same cases. We now turn our attention to the neutrinos, which are the main source of cooling. The results in this case are markedly different from what we initially found (Lee & Ramirez–Ruiz, 2002), simply because of the new equation of state, the more realistic cooling rates and the approximate computation of opacities. They are in general agreement with the preliminary results we presented before (Lee, Ramirez–Ruiz & Page, 2004), which used a less detailed equation of state than the one shown here. The optically thick region is present in every disk at the start of the calculation. As already mentioned, this has two important effects: enhancing the pressure and suppressing the neutrino emission. This reflects upon the total luminosity, since the suppression occurs precisely in the hottest regions, where most of the energy would otherwise be released. The most important qualitative difference between the runs presented here is that for a high viscosity ($`\alpha =0.1`$, runs a1M and a1m), the disk is drained of mass so fast that it has no chance to cool ($`t_{\mathrm{acc}}<t_{\mathrm{cool}}`$) and release most of its internal energy. It is in fact advected into the black hole. Moreover, the optically thick region disappears entirely from the disk (see Figure 5) by $`t=40`$ ms. The emission is then no longer suppressed (see the contours of $`\dot{q}`$ in the same figure), and the disk radiates at the maximum possible cooling rate. The thermal energy content of the disk is so large, that as it thins, the luminosity actually increases briefly around $`t=10`$ ms before dropping again. The drop at late times follows an approximate power law, with $`L_\nu t^1`$. In this interval the energy release comes from a combination of residual internal energy and viscous dissipation within the disk. For lower values of the viscosity, $`\alpha 0.01`$, the central regions of the disk remain optically thick throughout the calculations. Note that reducing the disk mass by a factor of five (as was done for runs a1m, a2m, a3m) does not affect these overall conclusions. The fluid is simply compressed into a smaller volume, and thus the densities and temperatures that are reached are similar than for the high mass runs. For run a3M the neutrino luminosity is practically constant at $`3\times 10^{52}`$ erg s<sup>-1</sup> for $`t100`$ ms. Table 2 summarizes the typical disk mass, energy density, accretion rate, luminosity, and duration and energetics of neutrino emission for all runs. ## 4 Summary, conclusions and astrophysical implications ### 4.1 Summary We have performed two-dimensional hydrodynamical simulations of accretion disks in the regime of hypercritical accretion, where neutrino emission is the main cooling agent. The disks are assumed to be present around a stellar–mass black hole, and are evolved for a few hundred dynamical timescales. We have paid particular attention to the relevant microphysical processes under the conditions at hand, and used a detailed equation of state which includes an ideal gas of $`\alpha `$ particles and free baryons in nuclear statistical equilibrium and a relativistic Fermi gas of arbitrary degeneracy. The composition of the fluid is determined by weak interactions. The density and temperature are such that the inner regions of the disk become opaque to neutrinos, and this is taken into account in a simple approximation. ### 4.2 Conclusions Our main conclusions can be summarized as follows: * Once the fluid becomes photodisintegrated into free nucleons, neutronization becomes important and lowers the electron fraction substantially below 1/2, with the electron fraction $`Y_e`$ reaching $`0.05`$ at its minimum. This value, however, does not occur in the immediate vicinity of the black hole, but rather at the transition radius where the fluid becomes optically thick to its own neutrino emission, $`r_{}10^7`$ cm. At smaller radii, the electron fraction rises again, reaching $`0.1`$ close to the horizon. * Neutrino trapping produces a change in composition and an inversion in the electron fraction. The associated negative gradient in $`Y_e(r)`$ induces convective motions in the optically thick region of the disk. This is analogous to what presumably occurs following core collapse, in a proto–neutron star and its surrounding envelope. Due to the radial flows induced by viscosity, convection is unable to suppress this composition gradient. It would appear, however, that the entropy per baryon in the optically thick region is very close to being constant, with $`s/k6`$. * The spatial structure of the disk is characterized by the transition radius $`r_{}`$ where the material becomes optically thick. For $`r>r_{}`$ the disk cools efficiently, whereas for $`r<r_{}`$ the emission is suppressed and the fluid is unable to cool directly (although it does so on a neutrino diffusion timescale). This leads to larger pressures and a more moderate rise in density. * The temporal evolution of the disk is determined by the balance between accretion and neutrino emission. For low viscosities ($`\alpha 0.01`$), the disk is able to cool in a quasi steady state and radiate its internal energy reservoir. This lasts for approximately 0.1-0.4 s, with $`L_\nu 10^{53}`$ erg s<sup>-1</sup>. Thereafter the typical luminosity and density quickly decay. For large viscosities ($`\alpha 0.1`$) the disk is drained of mass on an accretion timescale which is shorter than the cooling timescale, and the internal energy of the disk is essentially advected into the black hole. An interesting result in this case is that as the disk becomes transparent before being engulfed by the hole, it undergoes a re-brightening, as some of the stored energy escapes. * The total energy output in neutrinos is $`E_\nu 10^{52}`$ erg, over a timescale of $`200`$ ms. The typical accretion rates are $`0.1M_{\mathrm{}}`$ s<sup>-1</sup>, and neutrino energies are $`8`$ MeV at $`r_{}`$. Energy densities in the inner regions of the disk are $`10^{31}`$erg cm<sup>-3</sup>. ### 4.3 Discussion There are two main ingredients in the results presented here that contrast with those available previously in the literature concerning the steady state structure of neutrino cooled accretion disks. The first is the ability to dynamically model the evolution of the system for hundreds of dynamical timescales, taking the previous history of the fluid into consideration through the choice of initial conditions. This allows us to consider the energetics on more relevant timescales (cooling, viscous) than the dynamical one accessible in three dimensional studies. The second is the realization that the structure of the disks is affected qualitatively by the presence of an optically thick region at high densities. In this sense the situation is similar to that encountered following massive core collapse, where convection occurs (recent work assessing the relative importance of the MRI, neutrino and convection driven viscosity in rotating collapsing cores has been reported by Akiyama et al., 2003; Thompson, Quataert & Burrows, 2005, we comment further on this issue in §4.3.2). Clearly there is room for improvement in the results presented here. To begin with, our expressions for the effect of a finite optical depth on cooling and pressure are too simple, and attempt only to capture the essential physical behavior of the system. There is an important region of the disk where the optical depth is not large enough to consider the diffusion approximation, and where more detailed transport effects ought to be considered. We have not separated the neutrino variables into three species, which would be more rigorous (e.g., Yokosawa et al., 2004, have noted that in such flows, the electron neutrinos $`\nu _e`$ might be preferentially absorbed with respect to electron anti–neutrinos $`\overline{\nu }_e`$, thus affecting the neutrino annihilation luminosity above the surface of the disk in an important way). We have only considered the effects of coherent scattering off nucleons and $`\alpha `$ particles for opacity purposes. This specifically ignores absorptive scattering, which would: (i) lead to a modification of the neutrino emergent spectrum; (ii) produce heating of the fluid. The latter may be particularly important concerning the driving of powerful winds off the surface of the disk, and needs to be addressed more carefully in the context of GRBs. As in previous work, we have chosen to maintain the Newtonian expression for the potential well of the central black hole for ease of comparison to previous results (this will be addressed in future work). #### 4.3.1 Gamma-Ray Bursts The sudden release of gravitational binding energy of a neutron star is easily sufficient to power a GRB. The minimum energy requirement is $`10^{51}(\mathrm{\Delta }\mathrm{\Omega }/4\pi )`$ erg if the burst is beamed into a solid angle $`\mathrm{\Delta }\mathrm{\Omega }`$. As discussed in §1, a familiar possibility, especially for bursts belonging to the short duration category, is the merger of a neutron star - black hole, or double neutron star binary by emission of gravity waves, which, as illustrated here, is likely to generate a black hole surrounded by a lower mass accretion disk. How is the available rotational and gravitational energy converted into an outflowing relativistic plasma? A straightforward way is that some of the energy released as thermal neutrinos is reconverted, via collisions outside the disk, into electron-positron pairs or photons. The neutrino luminosity emitted when disk material accretes via viscous (or magnetic) torques on a timescale $`\mathrm{\Delta }t1`$ s is roughly $$L_\nu 2\times 10^{52}\left(\frac{M_{\mathrm{disk}}}{0.1M_{\mathrm{}}}\right)\left(\frac{\mathrm{\Delta }t}{1\mathrm{s}}\right)^1\text{erg s}\text{-1}$$ (30) for a canonical radiation efficiency of 0.1. During this time, the rate of mass supply to the central black hole is of course much greater than the Eddington rate. Although the gas photon opacities are large, the disk becomes sufficiently dense and hot to cool via neutrino emission. There is in principle no difficulty in dissipating the disk internal energy, but the problem is in allowing these neutrinos to escape from the inflowing gas. At sufficiently low accretion rates, $`\alpha 0.01`$, we find that the energy released by viscous dissipation is almost completely radiated away on a timescale given by $`t_{\mathrm{cool}}E_{\mathrm{int}}/L_\nu 0.1`$ s. In contrast, for a higher mass supply, $`\alpha 0.1`$, energy advection remains important until the entire disk becomes optically thin. The restriction on the cooling rate imposed by high optical depths is key because it allows the energy loss to be spread over an extended period of time during which the neutrino luminosity stays roughly constant. This gives a characteristic timescale for energy extraction and may be essential for determining the duration of neutrino-driven short GRBs (Lee, Ramirez–Ruiz & Page, 2004). Neutrinos could give rise to a relativistic pair-dominated wind if they converted into pairs in a region of low baryon density (e.g. along the rotation axis, away from the equatorial plane of the disk). The $`\nu \overline{\nu }\mathrm{e}^+\mathrm{e}^{}`$ process can tap the thermal energy of the torus produced by viscous dissipation. For this mechanism to be efficient, the neutrinos must escape before being advected into the hole; on the other hand, the efficiency of conversion into pairs (which scales with the square of the neutrino density) is low if the neutrino production is too gradual. Typical estimates suggest a lower bound<sup>3</sup><sup>3</sup>3This estimate, however, assumes that the entire surface area emits close to a single temperature black–body. It should be noted that if the dissipation takes place in a corona–like environment, the efficiency may be significantly larger (Ramirez-Ruiz & Socrates, 2005). of $`L_{\nu \overline{\nu }}10^3L_\nu `$ (e.g Rosswog & Ramirez–Ruiz, 2003). If the pair-dominated plasma were collimated into a solid angle $`\mathrm{\Delta }\mathrm{\Omega }`$ then of course the apparent isotropized energy would be larger by a factor $`(4\pi /\mathrm{\Delta }\mathrm{\Omega })`$, but unless $`\mathrm{\Delta }\mathrm{\Omega }`$ is $`10^2`$ this may fail to satisfy the apparent isotropized energy of $`10^{52}`$ ergs implied by a redshift $`z=1`$ for short GRBs. One attractive mechanism for extracting energy that could circumvent the above efficiency problem is a relativistic magneto hydrodynamic (MHD) wind (Usov, 1992; Thompson, 1994). Such a wind carries both bulk kinetic energy and ordered Poynting flux, and it possible that gamma-ray production occurs mainly at large distances from the source (Duncan & Thompson, 1992; Usov, 1994; Thompson, 1994). A rapidly rotating neutron star (or disk) releases energy via magnetic torques at the rate $`L_{\mathrm{mag}}10^{49}B_{15}^2P_3^4R_6^6\mathrm{erg}\mathrm{s}^1`$, where $`P=10^3P_3`$ s is the spin period, and $`B=10^{15}B_{15}`$ G is the strength of the poloidal field at a radius $`R=10^6R_6`$ cm. The last stable orbit for a Schwarzschild hole lies at a coordinate distance $`R=6GM/c^2=9(M/M_{\mathrm{}})`$ km, to be compared with $`R=GM/c^2=3/2(M/M_{\mathrm{}})`$ km for an extremal Kerr hole. Thus the massive neutron disk surrounding a Schwarzschild black hole of approximately $`2M_{\mathrm{}}`$ should emit a spin-down luminosity comparable to that emitted by a millisecond neutron star. A similar MHD outflow would result if angular momentum were extracted from a central Kerr hole via electromagnetic torques (Blandford & Znajek, 1977). The field required to produce $`L_{\mathrm{mag}}10^{51}\mathrm{erg}\mathrm{s}^1`$ is colossal, and may be provided by a helical dynamo operating in hot, convective nuclear matter with a millisecond period (Duncan & Thompson, 1992). A dipole field of the order of $`10^{15}`$ G appears weak compared to the strongest field that can in principle be generated by differential rotation ($`10^{17}[P/1\mathrm{ms}]^1\mathrm{G}`$), or by convection ($`10^{16}\mathrm{G}`$), although how this may come about in detail is not resolved. We examine in more detail the possible generation of strong magnetic fields below in §4.3.2. Computer simulations of compact object mergers and black hole formation can address the fate of the bulk of the matter, but there are some key questions that they cannot yet tackle. In particular, high resolution of the outer layers is required because even a tiny mass fraction of baryons loading down the outflow severely limits the attainable Lorentz factor - for instance a Poynting flux of $`10^{53}`$ erg could not accelerate an outflow to $`\mathrm{\Gamma }100`$ if it had to drag more than $`10^4`$ solar masses of baryons with it. One further effect renders the computational task of simulating jet formation even more challenging. This stems from the likelihood that the high neutrino fluxes ablate baryonic material from the surface of the disk at a rate (Qian & Woosley, 1996) $$\dot{M}_\eta 5\times 10^4\left(\frac{L_\nu }{10^{52}\mathrm{erg}\mathrm{s}^1}\right)^{5/3}M_{\mathrm{}}\mathrm{s}^1.$$ (31) A rest mass flux $`\dot{M}_\eta `$ limits the bulk Lorentz factor of the wind to $$\mathrm{\Gamma }_\eta =\frac{L_{\mathrm{mag}}}{\dot{M}_\eta c^2}=10\left(\frac{L_{\mathrm{mag}}}{10^{52}\mathrm{erg}\mathrm{s}^1}\right)\left(\frac{\dot{M}_\eta }{5\times 10^4M_{\mathrm{}}\mathrm{s}^1}\right)^1.$$ (32) If one assumes that the external poloidal field strength is limited by the vigor of the convective motions, then the spin-down luminosity scales with neutrino flux as $`L_{\mathrm{mag}}B^2v_{\mathrm{con}}^2L_\nu ^{2/3}`$, where $`v_{\mathrm{con}}`$ is the convective velocity. The ablation rate given in equation (31) then indicates that the limiting bulk Lorentz factor $`\mathrm{\Gamma }_\eta `$ of the wind decreases as $`L_\nu ^1`$. Thus the burst luminosity emitted by a magnetized neutrino cooled disk may be self-limiting. The mass loss would, however, be suppressed if the relativistic wind were collimated into a jet. This suggests that centrifugally driven mass loss will be heaviest in the outer parts of the disk, and that a detectable burst may be emitted only within a certain solid angle centered on the rotation axis (see e.g., Rosswog & Ramirez–Ruiz, 2003). #### 4.3.2 Generation of Strong Magnetic Fields There is also the question of magnetic fields, which we have not included, but should obviously be considered. The field in a standard disk is probably responsible for viscous stresses and dissipation, through the MRI. In this respect the current scenario should exhibit these characteristics. The MRI operates on an orbital timescale, and so the field would grow in a few tens of milliseconds. It may also be amplified by the convective motions described above, § 3.2. The saturation value for the field can be naively estimated as that at which its energy density is in equipartition with the gas, $`B^2/8\pi \rho c_s^2`$, or when the Alfven speed, $`v_A=B/\sqrt{4\pi \rho }`$ is comparable with the azimuthal velocity, $`v_\varphi `$. This gives $`B10^{16}`$ G. It is not clear at all, however, that the field amplitude will reach such high levels, because the magnetic Reynolds number is far beyond its critical value (where diffusion balances dynamo–driven growth), and amplification can lead to field expulsion from the convective region, thereby destroying the dynamo (Rädler et al., 2002). Precious little is known about the growth of magnetic fields at such overcritical levels (Reinhardt & Geppert, 2005), and a definitive answer will require the self–consistent inclusion of full MHD into the problem at hand (but with a level of resolution which may be well above present computational capabilities). It is not clear either that the magnetic shearing instability can generate a mean poloidal field as strong as $`B10^{16}`$ G, since to first order it does not amplify the total magnetic flux. The non-linear evolution of the instability depends sensitively on details of magnetic reconnection, and it has indeed been suggested that this can smooth reversals in the field on very small scales, pushing the dominant growing mode to much larger scales (Goodman & Xu, 1994). It is certainly possible, as shown here, that compact binary mergers do form a neutron disk that is hot enough to be optically thick to neutrinos, and convective instability is a direct consequence of the hot nuclear equation of state. A neutron disk is likely to be convective if the accretion luminosity exceeds $`10^{50}10^{51}\mathrm{erg}\mathrm{s}^1`$. Note that even if the accretion luminosity is lower, a hot, massive disk (such as those forming in collapsars, Woosley, 1993) would undergo a brief period of convection as a result of secular cooling (notice that convection is driven by secular neutrino cooling, whereas the MRI is powered by a release of shear kinetic energy). If the dense matter rotates roughly at the local Keplerian angular velocity, $`\mathrm{\Omega }(GM_{\mathrm{BH}}/r^3)^{1/2}`$, then $`L_{\mathrm{mag}}`$ is approximately independent of radius, and the required poloidal field for a given luminosity is $`B_{15}L_{\mathrm{mag},50}^{1/2}(M_{\mathrm{BH}}/M_{\mathrm{}})^1`$. If a period of convection is a necessary step in the formation of a strong, large-scale poloidal field, an acceptable model thus requires that the surrounding torus should not completely drain into the hole on too short of a timescale. Whether a torus of given mass survives clearly depends on its thickness and stratification, which in turn depends on internal viscous dissipation and neutrino cooling. A large amount of differential rotation (as may occur in newborn neutron stars or those in X–ray binaries, and is definitely the case in toroidal structures supported mainly by centrifugal forces), combined with short periods, may produce substantial magnetic field amplification (Kluźniak & Ruderman, 1998; Spruit, 1999). The energy transferred to the magnetic field is released in episodic outbursts when the buoyancy force allows the field to rise to the surface of the star or disk. The amplification of a magnetic field to such strong values would clearly have important consequences on the evolution and time variability of the disk and its energy output. It would probably lead to strong flaring and reconnection events accompanied by the release of large amounts of energy, if the growth time for field amplification, $`t_\mathrm{B}`$, is shorter than the accretion timescale, $`t_{\mathrm{acc}}`$ (otherwise the disk would drain of matter before the field had the chance to reach large values; in this case, the survival of a massive, rapidly rotating neutron star as the end-point of binary NS merger might be preferred over the prompt formation of a BH). An effective helical dynamo of the $`\alpha \mathrm{\Omega }`$ type should be favored by a low effective viscosity, $`\alpha `$, because, as stated in §3.2, the overturn time<sup>4</sup><sup>4</sup>4Recall that the resulting convective motions in a proto-neutron star are extremely vigorous, with an overturn time of $`1`$ ms. is $`t_{\mathrm{con}}20`$ms (this applies only if the disk is not fed with matter externally for a time longer than $`t_{\mathrm{acc}}`$, otherwise convection would also be able to amplify the magnetic field). #### 4.3.3 Nucleosynthesis Core collapse SNe and compact object mergers are natural astrophysical sites for the production of heavy elements (Lattimer & Schramm, 1974; Meyer & Brown, 1997; Freiburghaus et al., 1999; Rosswog et al., 1999; Lee, 2001). In particular, nucleosynthesis in neutrino–driven winds is an issue that may be relevant for iron–group elements, as well as for heavier nuclei through the r–process (Woosley & Hoffman, 1992). Initial investigations into this matter (Qian & Woosley, 1996) determined that the entropy in the outflow arising from a newborn neutron star was probably too low to give rise to the r–process efficiently. More recently, this problem has been addressed again in the specific case of collapsar or post–merger accretion disks, based on the results of analytical calculations of neutrino cooled disks in one dimension (Pruet, Thompson & Hoffman, 2004). If the wind consists of a uniform outflow driven from the surface of the disk, the entropy is too low, and essentially only iron–group elements are synthesized, in agreement with the earlier results. However, their results also indicate that in a bubble-type outflow (against a background steady wind), where episodic expulsion of material from the inner regions takes place, material may be ejected from the disk and preserve its low electron fraction, thus allowing the r-process to occur. The convective motions reported here, occurring in the optically thick portion of the disk, would represent one way such cells could be transported to the disk surface if they can move fast enough to preserve the neutron excess. Since the existence of a convection region is dependent upon the densities reached in the inner disk (so that it becomes opaque), the synthesized nuclei (iron group vs. heavier, r–process elements) could be a reflection of its absence/presence. We have benefited from many useful discussions and correspondence with U. Geppert, N. Itoh, K. Kohri, P. Kumar, A. MacFadyen, P. Mészáros, M. Prakash, M. Rees, T. Thompson, S. Woosley and A. Socrates. Financial support for this work was provided in part by CONACyT 36632E (WHL, DP) and by NASA through a Chandra Postdoctoral Fellowship award PF3-40028 (ERR). Part of this work was done during visits to the Institute for Advanced Study (WHL) and Instituto de Astronomía, UNAM (ERR), whose hospitality is gratefully acknowledged. We thank the anonymous referee for helpful comments and suggestions on the initial manuscript.
warning/0506/cond-mat0506200.html	ar5iv	text	# On the relationship between phase transitions and topological changes in one dimensional models ## I The models The one dimensional models we study are all defined by the Hamiltonian $`=_{i=1}^Np_i^2/2m+V(\{q\}_{i=1\mathrm{}N})`$ ($`m`$ is the mass of each particle), where $`V`$ is the potential energy. We consider two different classes of models. The first one, introduced by Burkhardt bur as a model for localization-delocalization transition of interfaces, is defined by the potential energy $`V^{(1)}`$: $$V^{(1)}(\{q\}_{i=1\mathrm{}N})=\underset{i=1}{\overset{N}{}}K\|q_{i+1}q_i\|+\underset{i=1}{\overset{N}{}}V_p^{(1)}(q_i),$$ (1) where $`K`$ measures the strength of the force between neighboring pairs, $`V_p^{(1)}(q)`$ is the on-site pinning potential, and periodic boundary conditions are assumed $`q_{N+1}q_1`$. We chose for $`V_p^{(1)}(q)`$ the following form $$V_p^{(1)}(q)=\{\begin{array}{cc}+\mathrm{}& \text{for }q0\hfill \\ 0& \text{for }0<q<L\hfill \\ U_0& \text{for }LqL+R\hfill \\ 0& \text{for }q>L+R\hfill \end{array}$$ (2) that generalizes the original form in Ref. bur introducing a parameter $`L`$ that gives the position of the pinning potential (a square well of depth $`U_0`$ and width $`R`$, see Figure 1) from the edge of the system. The case with $`L=0`$ coincides with the original Burkhardt confining model, while the non-confining case is retrieved in the $`L\mathrm{}`$ limit. The models of the second class are defined by the the potential energy $`V^{(2)}`$ and $`V^{(3)}`$ of the form: $$V^{(2,3)}(\{q\}_{i=1\mathrm{}N})=\underset{i=1}{\overset{N}{}}\frac{K}{2}(q_{i+1}q_i)^2+\underset{i=1}{\overset{N}{}}V_p^{(2,3)}(q_i).$$ (3) We consider two different versions of this model, one defined by the on-site Morse potential, introduced by Peyrard and Bishop as a simple model for DNA thermal denaturation pey ; pey1 (PB model) $$V_p^{(2)}(q)=U_o\{(e^{q/R}1)^21\};$$ (4) the other is a symmetric version of the former (SPB model) $$V_p^{(3)}(q)=U_o\{(e^{\|q\|/R}1)^21\},$$ (5) a slight modification of PB model that allows for a non-confined motion of the variables (see Fig. 2). We note that the introduction of the modulus in the Eq. 5 does not introduce discontinuities up to the second derivative of the potential. The quantities $`U_o`$ and $`R`$ determine respectively the energy and the length scales of the on-site potential (in the following all quantities will be reported in $`U_o`$ and $`R`$ units). We further chose $`m`$=1. In all the three cases a parameter of the Hamiltonian is related to the strength of the inter-particles interactions ($`K`$) and, for the case of $`V^{(1)}`$ a second parameter is the position of the pinning potential $`L`$. Specifically, the relevant quantity defining the relative weight of the on-site with respect to interparticles potentials is the dimensionless ratio $`\xi `$=$`KR/U_o`$ or $`\xi `$=$`KR^2/2U_o`$ for the potential models (1) or (2,3) respectively, while the position of the pinning potential is given by $`\zeta =L/R`$ for the potential model (1). The generalized Burkhardt model has been treated analytically, while the Peyrard-Bishop models are studied numerically. In the latter cases we performed isothermal molecular dynamics simulations using Nosé-Hoover thermostat at different temperatures for systems with $`N`$=$`500`$ degrees of freedom with periodic conditions $`q_{N+1}`$=$`q_1`$. We studied different values of the control parameter $`\xi `$, as an example the results are reported for $`\xi `$=$`0.05`$ and $`0.5`$, all the other $`\xi `$ values give results in qualitative agreement with the two reported examples. ## II Burkhardt model ### II.1 Thermodynamics The thermodynamics of the Burkhardt model is known since many years bur for both the $`\zeta =0`$ and $`\zeta =\mathrm{}`$ cases. The method of solution is briefly outlined below. The determination of the thermodynamic of systems described by a potential function of the form $$V(\{q\}_{i=1\mathrm{}N})=\underset{i=1}{\overset{N}{}}K\|q_{i+1}q_i\|+\underset{i=1}{\overset{N}{}}V_p(q_i),$$ (6) i. e. similar to the case (1) (Eq. 1), goes through the exploitation of the transfer matrix technique. Indeed, the configurational partition function $`𝒵`$ is given by: $$𝒵_N=𝑑q_1\mathrm{}𝑑q_Ne^{\beta V(\{q\}_{i=1\mathrm{}N})},$$ (7) that, defining the transfer ”matrix” $$𝒯(x,y)=e^{\beta K\|xy\|}e^{\beta [V_p(x)+V_p(y)]/2},$$ (8) can be written as: $$𝒵_N=𝑑q_1\mathrm{}𝑑q_N\underset{i=1}{\overset{N}{}}𝒯(q_i,q_{i+1}),$$ (9) recalling that $`q_{N+1}q_1`$. With this notation, the (configurational) free energy density $$f=\frac{1}{\beta N}\mathrm{log}(𝒵_N)$$ (10) in the thermodynamic limit is promptly written as $$f=\frac{1}{\beta }\mathrm{log}(\mathrm{max}\{\overline{\lambda }\})$$ (11) where $`\overline{\lambda }`$ is the set of eigenvalues of the transfer matrix, i. e. the eigenvalues of the integral equation $$𝑑y𝒯(x,y)\varphi (y)=\lambda \varphi (x).$$ (12) The latter equation, with the substitution $$\psi (x)=e^{\beta V_p(x)/2}\varphi (x),$$ (13) turns out to be $$𝑑ye^{\beta K\|xy\|}e^{\beta V_p(y)}\psi (y)=\lambda \psi (x).$$ (14) The next step is performed by noticing that the operator $`[d^2/dx^2+\beta ^2K^2]`$ applied to $`\mathrm{exp}(\beta K\|xy\|)`$ produces a delta-function: $$\left[\frac{d^2}{dx^2}+\beta ^2K^2\right]e^{\beta K\|xy\|}=2\beta K\delta (xy),$$ (15) thus by applying the previous operator to the integral equation 14, it can be transformed in a Schroedinger like differential equation: $$\left[\frac{d^2}{dx^2}+\beta ^2K^2\frac{2\beta K}{\lambda }e^{\beta V_p(x)}\right]\psi (x)=0.$$ (16) This equation must be solved with the conditions that i) the ”eigenfunction” $`\psi (x)`$ was normalizable, and, ii) the boundary condition (implicit in Eq. 14) $`\psi ^{}(0)/\psi (0)=\beta K`$ was fulfilled. In summary, the calculation of the thermodynamic of system defined by the potential energy of the form in Eq. 1 is reduced to the solution of a Schroedinger-like differential equation and, in particular, to the finding the largest eigenvalue of the original integral equation 14. In general, as the eigenvalues are continuous and smooth function of the parameters (among which the temperature), no phase transitions are expected unless the two largest among them cross each other. ### II.2 The $`\zeta =0`$ case. Let us now apply the procedure to the potential function in Eq. 1 for the case $`\zeta =0`$. We do not report the details of the calculation, as they are based on standard techniques for solving Schroedinger equation in Quantum Mechanics QM ; in summary the ”eigenvalues” $`\lambda `$ are determined by the equation: $$z(\lambda )=\beta K$$ (17) with $`z(\lambda )={\displaystyle \frac{f_1(P,Q)\mathrm{sin}(QR)f_2(P,Q))\mathrm{cos}(QR)}{f_3(P,Q))\mathrm{sin}(QR)+f_4(P,Q)\mathrm{cos}(QR)}},`$ (18) having defined $`f_1(P,Q)=Q^2`$ (19) $`f_2(P,Q)=PQ`$ $`f_3(P,Q)=P`$ $`f_4(P,Q)=Q`$ and $`Q(\lambda )=\sqrt{{\displaystyle \frac{2\beta K}{\lambda }}e^{\beta U_o}\beta ^2K^2}`$ (20) $`P(\lambda )=\sqrt{\beta ^2K^2{\displaystyle \frac{2\beta K}{\lambda }}}.`$ The only possibility for the function $`z(\lambda )`$ to be real (condition required for Eq. 17 to have solution) is that $`P`$ was real (if $`Q`$ become imaginary, $`z(\lambda )`$ is still real), thus it exists a solution to Eq. 17 only if $`P`$ is real. Therefore, when $`P`$ vanishes, the eigenvalues $`\lambda `$ disappear (more specifically, disappear the eigenvalues of the discrete spectrum, and only those of the continuum spectrum remain), and the (configurational) free energy is discontinuous. For each temperature, the condition $`P(\lambda )=0`$ is fulfilled for a ”critical” $`\lambda `$, given by: $$\lambda _c=\frac{2}{\beta K}.$$ (21) Thus, the equation for the largest eigenvalue at the ”critical” point is given by $`z(\lambda _c)=\beta _cK`$, or $$\sqrt{e^{\beta _cU_o}1}\mathrm{tan}\left[\beta _cKR\sqrt{e^{\beta _cU_o}1}\right]=1,$$ (22) which gives us the required equation for the critical (inverse) temperature $`\beta _c`$. This equation can be rearranged, introducing the control parameter $`\xi =KR/U_0`$, as: $$\xi =\frac{1}{\beta _cU_o}\frac{1}{\sqrt{e^{\beta _cU_o}1}}\mathrm{arctan}\left[\frac{1}{\sqrt{e^{\beta _cU_o}1}}\right].$$ (23) The plot of the critical temperature (in reduced units $`k_BT/U_o`$) as a function of $`\xi `$ is reported as full line in Fig. 3. The phase transition is of the localization-delocalization type. The particles, kept together by the $`K\|xy\|`$ term of the potential, for $`T<T_c`$ are pinned close to the square well, while, for $`T>T_c`$ are delocalized in the $`q`$-axis. A simple calculation leads to the value of the critical energy $`v_c`$ (the equilibrium energy $`v(T)`$ at the transition point $`v_c=v(T_c)`$). From Eq. 11, we have $$v(T)=\frac{(\beta f)}{\beta }=\frac{\lambda ^{}(\beta )}{\lambda (\beta )}$$ (24) where $`\lambda (\beta )`$ is the solution of Eq. 17. Close to the critical point, $`\lambda (\beta )=2/\beta K`$, thus $`\lambda ^{}(\beta )/\lambda (\beta )=1/\beta `$ and $$v_c=k_BT_c$$ (25) independently from the value of $`\xi `$. ### II.3 The $`\zeta 0`$ case. The calculation for the case of generic $`\zeta `$ values is quite similar to the previous one. Also in this case, the eigenvalues $`\lambda `$ are determined by an equation like Eq. 17, $`z(\lambda )=\beta K`$, with $`z(\lambda )`$ again given by Eq. 18 and with the $`f_n(P,Q)`$ functions ($`n`$=$`1,\mathrm{},4`$) given by $`f_1(P,Q)`$ $`=`$ $`P\left[Q^2\mathrm{cosh}(PR\zeta )P^2\mathrm{sinh}(PR\zeta )\right]`$ $`f_2(P,Q)`$ $`=`$ $`P^2Q\mathrm{exp}(PR\zeta )`$ $`f_3(P,Q)`$ $`=`$ $`\left[P^2\mathrm{cosh}(PR\zeta )Q^2\mathrm{sinh}(PR\zeta )\right]`$ $`f_4(P,Q)`$ $`=`$ $`PQ\mathrm{exp}(PR\zeta ).`$ (26) Obviously, Eqs. 26 recover Eqs. II.2 in the $`\zeta 0`$ limit. The same considerations on the reality of $`P(\lambda )`$ reported above apply to Eq. 26. Thus the condition $`P(\lambda )=0`$ define the critical value of the eigenvalue, $`\lambda _c=2/\beta K`$,and the equation for the critical temperature ($`z(\lambda _c)=\beta _cK`$) becomes: $`\sqrt{e^{\beta _cU_o}1}\mathrm{sin}\left(\beta _cKR\sqrt{e^{\beta _cU_o}1}\right)\times `$ (27) $`\times [\mathrm{cos}\left(\beta _cKR\sqrt{e^{\beta _cU_o}1}\right)`$ $`\beta _cKR\zeta \sqrt{e^{\beta _cU_o}1}\mathrm{sin}\left(\beta _cKR\sqrt{e^{\beta _cU_o}1}\right)]^1=1.`$ Similar to the $`\zeta =0`$ case, this equation can be rearranged, introducing the control parameter $`\xi `$, as: $$\xi =\frac{1}{\beta _cU_o}\frac{1}{\sqrt{e^{\beta _cU_o}1}}\mathrm{arctan}\left(\frac{1}{\sqrt{e^{\beta _cU_o}1}}\frac{1}{1+\beta _cU_o\xi \zeta }\right).$$ (28) At variance with Eq. 23, this equation cannot be cast in the form $`\xi =\xi (\beta _c)`$, thus it must be solved numerically to plot the critical temperature as a function of the control parameter $`\xi `$. This plot is reported in Fig. 4 for different values of $`\zeta `$. As can be observed in Fig. 4, on increasing $`\zeta `$, i. e. on displacing the position of the square well towards high value of the coordinate, the critical temperature, for a given $`\xi `$ value, increases, expanding the amplitude of the ”cold” (localized or pinned) phase. This can be better seen in Fig. 5, where the $`\zeta `$ dependence of the critical temperature is reported for some values of $`\xi `$. We conclude this section noticing that the phase transition actually exists for all the value of $`\zeta `$, and in the limit of $`\zeta \mathrm{}`$, the critical temperature goes without discontinuities to infinity. Therefore, we are lead to conclude that the model investigated in Ref. kastner to demonstrate the unattainability of a purely topological criterion for the existence of a phase transition is a “borderline” model, in which the phase transition can be thought to be present at “$`T`$ infinity” (even though the precise meaning of this statement is not well defined). Then, to the same topology (as we will see in the next section, the topology of the potential function in Eq. 2 does not depend on the value of $`\zeta `$) always corresponds a phase transition. To discuss the question of the coincidence (or not) of the critical energy with the topological discontinuity, we need to calculate $`v_c(\xi ,\zeta )`$. Following the same argument reported for the case $`\zeta =0`$ we conclude that the critical potential energy depends on $`\xi `$ and $`\zeta `$ only through $`T_c`$: $`v_c=k_BT_c`$. As an example, in Fig. 6 we report the caloric curve $`v(T)`$ as a function of the inverse temperature for different $`\zeta `$ value and for $`\xi =1`$. For all the $`\zeta `$ values, on the low-$`\beta `$ side the curves end at the points ($`\beta _c`$, $`v_c`$); these points are aligned along the $`v_c(\beta _c)`$ line (thick dotted line) given by $`v_c(\beta _c)=1/\beta _c`$. ### II.4 Topology The analysis of the topological properties of the Burkhardt model is reported by Kastner in Ref. kastner . He analyzed only the two limiting cases of confining and non-confining models, corresponding in our notation to $`\zeta `$=$`0`$ and $`\zeta `$=$`\mathrm{}`$ respectively. He found that a topology change is present in both cases, even if not really equal in “strength”. The value of the potential energy at which the topological change appears is $`v_\theta `$=$`0`$, irrespective of the considered model. One can easy generalize the above analysis to the general case with arbitrary $`\zeta `$, and conclude that the topological change is always located at energy $`v_\theta `$=$`0`$. It is worth noting that the energy at which topological change appears is lower than the thermodynamic transition energy: $`v_c>v_\theta `$ (see Fig. 6). We will further discuss this issue in Sec. V, after having described the thermodynamics and topology of the PB and SPB models. ## III Peyrard-Bishop model ### III.1 Thermodynamics The thermodynamics of the Peyrad-Bishop model (defined in Eq.s 3, 4) can be studied using transfer matrix techniques, as the Burkhardt model described in the previous section. However, in this case approximated methods have to be considered in order to obtain a corresponding Schroedinger like differential equation. In the region $`\xi 1`$ and temperature window $`U_ok_BT\xi U_o`$ the classical statistical mechanics problem is mapped to the quantum Morse oscillator problem pey2 ; theo . Similarly to the case of Burkhardt potential, the presence of a second order phase transition for the Peyrad-Bishop model is signaled by the bounded-unbounded transition of the lower state in the corresponding quantum problem. In the above range of $`\xi `$ and $`T`$, Peyrard and Bishop obtained an analytical expression for the transition temperature $`k_BT_c/U_o`$=$`4\sqrt{\xi }`$ and transition energy $`v_c/U_o`$=$`k_BT_c/2U_o`$=$`2\sqrt{\xi }`$. For generic $`(\xi ,T)`$ values, only numerical results can be used to infer the existence and location of a phase transition. In Fig. 7 we report the temperature dependence of the potential energy per particle $`v=V/N`$ (full symbols) of the PB model for two different values of $`\xi `$: $`0.05`$ (upper panel) and $`0.5`$ (lower panel). Also reported in the figure are the energy $`v_c`$ (dot-dashed line) and temperature $`T_c`$ (full line) of the phase transition point: $`v_c/U_o0.61`$ and $`k_BT_c/U_o1.22`$ for $`\xi `$=$`0.05`$, $`v_c/U_o1.59`$ and $`k_BT_c/U_o3.20`$ for $`\xi `$=$`0.5`$. Dashed lines are the $`T`$-dependence of the potential energy in the high $`T`$ phase: $`v(T)`$=$`k_BT/2`$. ### III.2 Topology The topology of the Peyrard-Bishop model is studied in the paper of Grinza and Mossa gri\_mos . A topological change is found at the energy value $`v_\theta `$=$`0`$, corresponding to a topological change in the hypersurfaces $`\mathrm{\Sigma }_v`$ varying $`v`$: from a close hypersurface for $`v<v_\theta `$ to an open one for $`vv_\theta `$ gri\_mos . In Fig. 7 the value of $`v_\theta `$ is indicated by an horizontal dotted line. We note that, also in this case, the topological discontinuity is lower in energy than the thermodynamic one: $`v_c>v_\theta `$. ## IV Symmetric Peyrard-Bishop model ### IV.1 Thermodynamics The Symmetric Peyrard-Bishop model defined by Eq.s 3, 5 does not exhibit phase transition at finite $`T`$. This can be view from the fact that there is always a bound state in the corresponding quantum problem, in analogy with the non-confined Burkhardt model bur ; cuesta . In Fig. 8 the same quantities as in the PB case are reported for the SPB model: energy $`v`$ (full symbols) for $`\xi `$=$`0.05`$ (upper panel) and $`\xi `$=$`0.5`$ (lower panel). It is evident in this case the absence of a phase transition at finite $`T`$ (in the $`T`$-range investigated). ### IV.2 Topology Following a similar argument as in Ref. kastner , one can see that also in the SPB case one has a topological change at exactly the same energy level as in the PB model $`v_\theta `$=$`0`$ (even if not identical in strength to the previous one). We refer to the papers in Ref.s kastner ; gri\_mos for a more detailed discussion of the topology. In Fig. 8 the value of $`v_\theta `$ is indicated by an horizontal dotted line. ## V Underlying saddles In this section we study the properties of the stationary points visited by the systems. The concept of “underlying saddles” was first introduced in the study of glassy disordered systems sad\_lj ; sad\_cav ; cavagna to better understand the topological counterpart of the dynamic transition taking place in these systems. Recently, it has been applied also in the analysis of models that exhibit thermodynamic phase transitions, in order to emphasize the role of topological changes at the “underlying saddles” energy in driving the phase transition ktrig ; ktrig2 ; phi4 . Here we apply the same methodology to investigate the one dimensional systems introduced before. Let start with the models having a continuous potential energy function, the PB and SPB models, which allow for the usual definition of stationary points. At the end of the section we will extend the argument to the discontinuous case of Burkhardt model. ### V.1 Peyrard-Bishop and Symmetric Peyrard-Bishop models There are only two stationary points in the potential energy hypersurface of both models: a minimum located at $`q_1`$=$`q_2`$=$`\mathrm{}`$=$`q_N`$=$`0`$ and a saddle (with degenerate Hessian matrix) at $`q_1`$=$`q_2`$=$`\mathrm{}`$=$`q_N`$=$`\mathrm{}`$ gri\_mos . In order to associate one of the two stationary points to each instantaneous configuration of the system, we used a similar trick as in the analysis of glassy systems sad\_lj or mean-field models ktrig ; ktrig2 ; phi4 . In the latter one minimized the pseudo-potential $`W`$=$`\|V\|^2`$ during the dynamic evolution at different temperatures, so introducing a map from equilibrium energy levels to saddles energy levels: $`:v(v)v_s`$. Due to the peculiarity of the present models, where the saddle point is “infinitely” far from each equilibrium configuration, we decided to apply the $`W`$ minimization method in a two steps procedure: i) first we minimized the $`W_{int}`$ quantity defined using the interaction potential part of $`V`$, $`W_{int}`$=$`\|V_{int}\|^2`$, where $`V_{int}`$=$`_{i=1}^N\frac{K}{2}(q_{i+1}q_i)^2`$; ii) then we minimized the $`W_p`$ defined using the on site potential $`W_p`$=$`\|V_p^{(2,3)}\|^2`$. This procedure ensures that the point reached is a true stationary point, i.e. the minimum or the saddle. Obviously, this is a quite arbitrary definition of basins of attraction of stationary points. As said in the introduction, the robustness of the results with respect to the possible choices of definition of a saddle basin of attraction is still an open problem. In Fig. 7 the temperature dependence of the energy $`v_s`$ (open symbols) of underlying saddles is shown for the case $`\xi `$=$`0.05`$ (upper panel) and $`\xi `$=$`0.5`$ (lower panel) in the PB model. The remarkable fact is that at $`T_c`$ (vertical full line in Fig. 7) the identity $`v_s`$=$`v_\theta `$ holds. The map $`(v)`$ is shown for PB model (open symbols) in Fig. 9 for the two cases $`\xi `$=$`0.05`$ (upper panel) and $`\xi `$=$`0.5`$ (lower panel). One observe that, as before pointed out, one has $`(v_c)`$=$`v_\theta `$ for both $`\xi `$ values. The fact that $`v_s(T)`$ in Fig. 7, as well as $`(v)`$ in Fig. 9, has a “smooth” transition between its low $`T`$ (or $`v`$) and high $`T`$ (high $`v`$) regions is most likely due to a finite size effect ($`N`$=$`500`$ here) and both $`v_s(T)`$ and $`(v)`$ will probably tend towards a step function in the thermodynamic limit. The previous finding indicates that the relevant quantity to consider when we are looking for topological changes related to a phase transition is the underlying stationary point energy, obtained trough a map from the critical level $`v_c`$. It is worth noting that the map $``$ is constant ($`(v)`$=$`v_\theta `$) for a broad range of values, also below $`v_c`$, at variance with other cases where around the transition point the properties of visited saddles change XY ; ktrig ; ktrig2 ; phi4 . One can conjecture that the flatness of $`(v)`$ is a pathology of these one-dimensional models, that have a number of stationary points that is not extensive in $`N`$ (actually there are only $`2`$ stationary points). In Fig 8 we report the same quantities $`v_s`$ as before (open symbols), now for the SPB model, with $`\xi `$=$`0.05`$ (upper panel) and $`\xi `$=$`0.5`$ (lower panel). In this case no phase transition is present, and indeed the topological singularity is never visited, $`v_s(T)<v_\theta `$ for each finite temperature ($`T<\mathrm{}`$). ### V.2 Burkhardt model To apply the analysis of the previous section also to the Burkhardt model, one has to find a suitable definition of “saddles” and of “basin of attraction of a saddle” for a discontinuous potential. One possibility is the following: we first minimize the interaction potential $`V_{int}=_{i=1}^NK\|q_{i+1}q_i\|`$, which is equivalent to put all the $`q_i`$ equal to the center of mass coordinate $`\overline{q}=N^1_iq_i`$. If $`\overline{q}`$ lies in the well of the potential, i.e. $`\overline{q}[L,L+R]`$, we will associate the “minimum” to the initial configuration, otherwise we will associate it to the “saddle” (we use this terminology by analogy with the PB model). It is clear that the average energy of the “underlying saddles” is simply the average of the on-site energy of the center of mass coordinate, $$v_s(T)=V_p^{(1)}(\overline{q})_T.$$ (29) In the thermodynamic limit the center of mass $`\overline{q}`$ is peaked around its mean value and then we can substitute the right hand side of Eq. 29 with $`V_p^{(1)}(\overline{q})`$, a quantity that can be explicitly computed. To determine $`\overline{q}`$ we can use the distribution probability $`\|\varphi (x)\|^2`$, where $`\varphi (x)`$=$`e^{\beta V_p(x)/2}\psi (x)`$ and $`\psi (x)`$ is the eigenfunction of the transfer matrix operator corresponding to the maximum eigenvalue bur (see Sec. II A). We note that the saddle energy $`v_s(T)`$ is a step function, equals to the minimum energy $`U_o`$ when $`\overline{q}`$ lies inside the square well and equals to the saddle energy $`0`$ otherwise. The temperature $`T_J`$ at which the visited “underlying saddle” jumps from minimum to saddle is shown in Fig. 3 (dashed line) as a function of the parameter $`\xi `$ for the $`\zeta `$=$`0`$ case. It is worth noting that the temperature $`T_J`$ lies always below the thermodynamic transition temperature $`T_c`$ (in analogy with the PB model, see Fig. 7). The same happens for all values of $`\xi `$. Therefore, also for the Burkhardt case, at the transition temperature $`T_c`$ the “underlying saddles” lie at an energy equal to the topological discontinuity energy $`v_\theta `$, i.e. $`(v_c)`$=$`v_\theta `$. ## VI Conclusions Studying two particular one dimensional models discussed in the recent literature kastner ; gri\_mos (Burkhardt model in the confining and non-confining version, Peyrard-Bishop model and its non-confining counterpart), we have focused on the relationship between phase transitions and topological changes, recently proposed in the literature cccp ; fps ; cpc ; fra\_pet . In these models, a topological singularity at a given energy value $`v_\theta `$(=0) is always found; however, i) in the confining version a phase transition is found but the critical energy is $`v_c>v_\theta `$ gri\_mos ; ii) in their non-confining version there is no phase transition at any finite temperature kastner . These results generated confusion as i) was interpreted as a confirmation of the strong topological hypothesis of Pettini et al. nota while ii) was considered as an evidence for the unattainability of a purely topological criterion for detecting phase transitions, although demonstrated only for the particular non-confining one dimensional models. Exploiting the concept of “underlying stationary points” defined through a generalization of the methods used in the glassy literature (minimization of the pseudopotential $`W`$=$`\|V\|^2`$), we have defined a map $`:vv_s`$ from energy level $`v`$ of $`V`$ to stationary points, with energy $`v_s`$. We have shown that: i) in the confining case, where the phase transition is present, one has $`(v_c)=v_\theta `$, in agreement with the weak topological hypothesis; ii) in the non-confining case, where the phase transition is not present at finite temperature (as the transition temperature goes continuously to infinity when the confining wall is removed) the energy of the underlying saddles is always below the topological singularity, i.e. $`v_s(T)<v_\theta ,T`$; the singular point $`v_\theta `$ is indeed visited for $`T\mathrm{}`$, consistently with the observation that the critical temperature is “infinite” in the non-confining case. The weak topological hypothesis appears as a possible framework to fit the results that recently appeared in the literature on all the different models investigated so far. Within this hypothesis three different scenarios are possible: 1. If there is no topological singularity $`v_\theta `$, a phase transition is not possible; this is consistent with the hypothesis of Pettini et al.: topological singularities are necessary conditions for a phase transition to take place. 2. If there is a topological singularity at energy $`v_\theta `$, a phase transition is also present if and only if there exist a temperature $`T_c`$ such that $`v_s(T_c)=v_\theta `$ (or equivalently an energy $`v_c`$ such that $`(v_c)=v_\theta `$). The above findings seem to indicate that, at least for the particular models investigated, a sufficiency criterion for the phase transition to take place requires the introduction of a statistical measure: thus, we believe that the statement of Kastner kastner concerning the unattainability of a purely topological criterion for detecting phase transitions is indeed correct, even though in Ref. kastner it has been derived using a “borderline” model (see Section II C). Let us conclude with two remarks: i) as already stated, the definition of the map $``$ is not unique, different definitions giving (slightly) different results. Thus, the weak topological hypothesis contains in its formulation an ambiguity and must be regarded only as a practical tool, at least at this stage of comprehension; ii) nevertheless, we hope that this approach can be of interest for the numerical investigation of systems of “mesoscopic” size (e.g. proteins and large molecules), i.e. such that the number of degrees of freedom is not large enough to allow to detect the presence of a phase transition using standard techniques. We thank M. Pettini and M. Kastner for helpful comments and suggestions.
warning/0506/cond-mat0506376.html	ar5iv	text	# Excitonic effects on the two-color coherent control of interband transitions in bulk semiconductors ## I Introduction The phenomenon of quantum interference can be used to control physical and chemical processes. One approach, the ‘$`n+m`$’ scheme, uses a two-color light field to interfere $`n`$\- and $`m`$-photon transitions Manykin and Afanasev (1967); Shapiro and Brumer (1997); Gordon et al. (1999). Interference between one- and two-photon transitions, for example, allows controllable polar asymmetry of photoelectrons in atomic ionization Yin et al. (1992); Baranova et al. (1992), controllable dissociation of HD<sup>+</sup> Sheehy et al. (1995), and controllable photocurrent injection in unbiased solids due to free carrier absorption Baskin and Éntin (1988), impurity-band transitions Éntin (1989), and quantum well bound-continuum intersubband transitions Dupont et al. (1995). In biased asymmetric semiconductor double wells, ‘$`1+2`$’ interference allows control of carrier population and THz emission Pötz (1998). Interband ‘$`1+2`$’ interference in unbiased semiconductors, which is our interest here, allows independent control of electrical current injection Atanasov et al. (1996); Haché et al. (1997) and spin current injection Bhat and Sipe (2000); Stevens et al. (2002); Stevens et al. (2003a); Hübner et al. (2003); Najmaie et al. (2003); Marti et al. (2004). Furthermore, in noncentrosymmetric semiconductors, it allows independent control of carrier populations (i.e., absorption) Fraser et al. (1999); Fraser and van Driel (2003) and carrier spin polarization Stevens et al. (2003b); Stevens et al. (2004). In each scenario, the experimenter can control the interference by adjusting the phases of the two colors. The controllable optical phases are not the only source of phase between the transition amplitudes. In general, there is also a material-dependent intrinsic phase Shapiro et al. (1988). Phenomenologically, the intrinsic phase appears as a phase shift in the dependence of the process on a relative phase parameter of the optical fields. Additionally, selectivity between two processes is possible when their intrinsic phases differ Shapiro and Brumer (2003); for example, ‘1+3’ experiments on diatomic molecules have controlled the branching ratio of ionization and dissociation channels Zhu et al. (1995). The intrinsic phase can be strongly frequency dependent near resonances Zhu et al. (1995), and the hope that it might serve as a new spectroscopic observable Seideman (1998, 1999) has led to efforts to understand its physical origin. Whereas a resonance is necessary for a phase shift to a ‘1+3’ process Seideman (1998), it is not necessary for a phase shift to a ‘1+2’ process. For example, an intrinsic phase in the ‘1+2’ photoionization of atoms is predicted from the simple model of a delta function potential Baranova et al. (1990); Pazdzersky and Usachenko (1997). Nevertheless, microscopic theories for the interband ‘$`1+2`$’ processes in bulk semiconductors have until now predicted trivial intrinsic phases Atanasov et al. (1996); Fraser et al. (1999); Bhat and Sipe (2000); Stevens et al. (2003b). The photocurrent, for example, was predicted to be proportional to the sine of the relative phase parameter for all final energies Atanasov et al. (1996). Each of these theories use the independent particle approximation, in which the Coulomb attraction between the injected electron and hole is neglected. That approximation is expected to be good for final energies well above the band gap, since in this case the electron and hole travel quickly away from each other. However, close to the band gap, one generally expects to see deviations from the independent particle approximation. In the one-photon absorption spectrum, for instance, it is well known that the electron-hole attraction is responsible for exciton lines below the band gap, and for an enhancement of the absorption above the gap known as Sommerfeld or Coulomb enhancement Haug and Koch (1993). The effect of the electron-hole attraction on one-photon absorption has been studied with various degrees of sophistication. On the one hand, modern ab initio calculations that include Coulomb effects have recently given very good quantitative agreement with experimental spectra Albrecht et al. (1998); Benedict et al. (1998a, b); Rohlfing and Louie (1998, 2000), although at the cost of significant computational overhead. On the other hand, simple models of Wannier excitons can describe Coulomb effects near the band edge of many direct gap semiconductors. These excitonic effects have long been understood qualitatively on the basis of a simple two-band model in the effective mass approximation Elliott (1957), which is even quantitatively accurate for typical semiconductors Rohlfing and Louie (1998). Excitonic effects on nonlinear optical properties of bulk semiconductors have also been studied in the effective mass, Wannier exciton approximation Loudon (1962); Mahan (1968); Rustagi et al. (1973); Lee and Fan (1974); Doni et al. (1974); Sondergeld (1977a); Girlanda et al. (1981); Blossey (1970); Ganguly and Birman (1967); Martin (1971); García-Cristóbal et al. (1994, 1998); Kolber and Dow (1978); Sheik-Bahae et al. (1994) and only recently with ab initio methods Chang et al. (2002). The two-photon absorption spectrum shows a different set of exciton lines and a Coulomb enhancement that is weaker than its one-photon counterpart. One- and two-photon absorption spectra have been measured sufficiently often that excitonic effects on them are well established. In contrast, semiconductor ‘1+2’ interference experiments have been done typically at only a single energy and typically many exciton binding energies above the band-gap. Moreover, these initial experiments lacked an absolute calibration of the relative optical phase, and thus were insensitive to the intrinsic phase. Such a calibration is possible Schumacher and Bucksbaum (1996), and could be used to verify the predictions we present here. A nontrivial intrinsic phase would have implications for the use of ‘$`1+2`$’ current injection as an absolute measurement of the carrier-envelope phase of an ultrashort optical pulse Roos et al. (2003); Fortier et al. (2004); Roos et al. (2005). In the present work, we extend the theory of ‘1+2’ coherent control of bulk semiconductor interband transitions beyond the independent particle approximation, employing a set of approximations that are valid close to the $`\mathrm{\Gamma }`$ point of a direct gap bulk semiconductor. Our investigation is limited to a perturbative treatment in the fields. In this limit of low photoinjected carrier density, the only inter-particle interaction of importance is that between a single electron and hole. We show that, due to the electron-hole attraction, a nontrivial phase shift does in fact occur in the control of current and spin current, but not in the control of carrier population or spin polarization. The intrinsic phase can be understood in terms of partial wave phase shifts due to the Coulomb potential between electron and hole. In addition, we find an enhancement of each process, and relate it to the Coulomb enhancements of one- and two-photon absorption. In the next section, we establish notation necessary to describe the ‘1+2’ processes in terms of one- and two-photon transition amplitudes. In section III, we present the microscopic model. The transition amplitudes are worked out in section IV. The final expressions for the ‘1+2’ effects are given in section V, and numerical results for GaAs are presented. In section VI, further understanding of the enhancement and intrinsic phase is discussed, and we examine the ratios often used as figures of merit for ‘1+2’ effects. Intermediate state Coulomb enhancement is examined in Appendix A. For reference, the current injection tensor is worked out for parabolic bands in Appendix B. ## II Preliminaries The rate of photocurrent injection into an unbiased bulk semiconductor by a two color light field $`𝐄(t)=𝐄_\omega \mathrm{exp}(i\omega t)+𝐄_{2\omega }\mathrm{exp}(i2\omega t)+c.c.`$ can be written $$\frac{dJ}{dt}^i=\eta _{(I)}^{ijkl}E_\omega ^jE_\omega ^kE_{2\omega }^l+c.c.,$$ (1) where superscripts denote Cartesian components and repeated indices are to be summed over Atanasov et al. (1996). The fourth rank tensor $`\eta _{(I)}`$, called the current injection tensor, describes the material response. It is purely imaginary in the independent particle approximation Atanasov et al. (1996), but can be complex in general. We define the intrinsic phase $`\delta ^{ijkl}`$ of the component $`\eta ^{ijkl}`$ as $$\delta ^{ijkl}=\mathrm{arctan}\left(\mathrm{Re}\left(\eta _{(I)}^{ijkl}\right)/\mathrm{Im}\left(\eta _{(I)}^{ijkl}\right)\right)$$ (2) so that it is zero or $`\pi `$ in the independent particle approximation. When the electron-hole interaction is included in the set of approximations used below, all the components of $`\eta _{(I)}`$ have the same phase. That is, $$\eta _{(I)}^{ijkl}=ie^{i\delta }\left\|\eta _{(I)}^{ijkl}\right\|.$$ (3) The intrinsic phase $`\delta `$ appears as a phase shift in the dependence of the current injection on the phase of the optical fields. For co-linearly polarized fields ($`𝐄_\omega =E_\omega \mathrm{exp}(i\varphi _\omega )\widehat{𝐱}`$ and $`𝐄_{2\omega }=E_{2\omega }\mathrm{exp}(i\varphi _{2\omega })\widehat{𝐱}`$), for example, the current injection is $$\frac{d𝐉}{dt}=2E_\omega ^2E_{2\omega }\left\|\eta _{(I)}^{xxxx}\right\|\mathrm{sin}\left(2\varphi _\omega \varphi _{2\omega }\delta \right)\widehat{𝐱}.$$ We are ignoring scattering processes through which the current would relax to a steady state value under continuous illumination, or would decay to zero following pulsed excitation. The current injection discussed here can be used as a source term in hydrodynamic equations that treat the scattering phenomenologically Atanasov et al. (1996); Haché et al. (1998); Côté et al. (1999) or in microscopic transport equations Král and Sipe (2000). Coulomb effects other than the excitonic effects we consider here play a role in scattering, especially at high densities of excited carriers. Such Coulomb effects are outside the scope of this paper. The current injection tensor can be written in terms of one- and two-photon transition amplitudes. We take as the initial state a clean, cold semiconductor. In a Fermi’s golden rule calculation for the ballistic current, light produces transitions to final states $`\|n`$ with velocity $`𝐯_{nn}`$, probability amplitude $`c_n`$, and energy $`\mathrm{}\omega _n`$ above the ground state. The final states will be specified in detail in the next section; here the label $`n`$ represents the set of quantum numbers for either an interacting or independent electron-hole pair. Thus, $`{\displaystyle \frac{d}{dt}}𝐉`$ $`=`$ $`{\displaystyle \frac{e}{L^3}}{\displaystyle \underset{n}{}}𝐯_{nn}{\displaystyle \frac{d}{dt}}\left\|c_n\right\|^2`$ $`=`$ $`{\displaystyle \frac{2\pi e}{L^3}}{\displaystyle \underset{n}{}}𝐯_{nn}\left\|\mathrm{\Omega }_n^{(1)}+\mathrm{\Omega }_n^{(2)}\right\|^2\delta \left(2\omega \omega _n\right)`$ $`=`$ $`{\displaystyle \frac{2\pi e}{L^3}}{\displaystyle \underset{n}{}}𝐯_{nn}\left\{\right\|\mathrm{\Omega }_n^{(1)}\|^2+\|\mathrm{\Omega }_n^{(2)}\|^2+(\mathrm{\Omega }_n^{(1)}\mathrm{\Omega }_n^{(2)}+c.c.)\}\delta (2\omega \omega _n),`$ where $`e`$ is the electron charge (negative), $`L^3`$ is a normalization volume, and $`\mathrm{\Omega }_n^{\left(i\right)}`$ is the amplitude for an i-photon transition. The $`\mathrm{\Omega }_n^{\left(i\right)}`$ take the form $`\mathrm{\Omega }_n^{(1)}`$ $`=`$ $`𝐄_{2\omega }𝐃_n^{\left(1\right)}`$ (5) $`\mathrm{\Omega }_n^{(2)}`$ $`=`$ $`𝐄_\omega 𝐄_\omega :𝖣_n^{\left(2\right)},`$ (6) where the vector $`𝐃_n^{\left(1\right)}`$ and second rank tensor $`𝖣_n^{\left(2\right)}`$ depend only on properties of the material. We consider them in detail in section IV. In the last equation of (II), the first term is (one-photon) one-color current injection, the second term describes two-photon one-color current injection, while the interference term describes the ‘1+2’ current. Because of their different dependencies on the electric field amplitudes, these three terms can in principle be separated experimentally. But the first two terms vanish for centrosymmetric crystals, and the first vanishes even for zincblende crystals; we neglect them here. Excitonic effects on the first term were studied by Shelest and Éntin Shelest and Éntin (1979); Belinicher and Sturman (1980). The third term survives in all materials. Comparing its expression with the phenomenological form (1), we have $$\eta _{(I)}^{ijkl}=\frac{2\pi e}{L^3}\underset{n}{}v_{nn}^i\left(D_n^{(2)}\right)^{jk}\left(D_n^{(1)}\right)^l\delta \left(2\omega \omega _n\right).$$ (7) Even if scattering from impurities and phonons is neglected, the injection current described by (II) does not capture the full current density; there are also optical rectification and ‘shift’ contributions to the current Aversa and Sipe (1996). The one-color varieties of these have been studied in some detail Belinicher and Sturman (1980); Sturman and Fridkin (1992); Sipe and Shkrebtii (2000); Côté et al. (2002), but the ‘1+2’ varieties have not. However, the different time dependencies of the three current contributions allows for their separate examination experimentally, at least in principle, and rough order-of-magnitude estimates indicate that typically the injection current will be the largest; it is the only contribution to the current we treat here. As in (II), the carrier injection rate can be written as $`{\displaystyle \frac{d}{dt}}n`$ $`=`$ $`{\displaystyle \frac{1}{L^3}}{\displaystyle \underset{n}{}}{\displaystyle \frac{d}{dt}}\left\|c_n\right\|^2`$ $`=`$ $`{\displaystyle \frac{2\pi }{L^3}}{\displaystyle \underset{n}{}}\left\{\right\|\mathrm{\Omega }_n^{(1)}\|^2+\|\mathrm{\Omega }_n^{(2)}\|^2+(\mathrm{\Omega }_n^{(1)}\mathrm{\Omega }_n^{(2)}+c.c.)\}\delta (2\omega \omega _n),`$ where $`n`$ is the number density of injected electron-hole pairs. The first two terms in (II) are the usual one- and two-photon absorption rates, which we denote by $`\dot{n}_{2\omega }`$ and $`\dot{n}_\omega `$ respectively. In terms of one- and two-photon coefficients $`\xi _{\left(1\right)}^{ij}`$ and $`\xi _{\left(2\right)}^{ijkl}`$ that depend only on the properties of the material, they can be written as $`\dot{n}_{2\omega }=\xi _{\left(1\right)}^{ij}E_{2\omega }^iE_{2\omega }^j`$ and $`\dot{n}_\omega =\xi _{\left(2\right)}^{ijkl}E_\omega ^iE_\omega ^jE_\omega ^kE_\omega ^l`$ Atanasov et al. (1996). From (II), $`\xi _{\left(1\right)}^{ij}`$ $`=`$ $`{\displaystyle \frac{2\pi }{L^3}}{\displaystyle \underset{n}{}}\left(D_n^{(1)}\right)^i\left(D_n^{(1)}\right)^j\delta \left(2\omega \omega _n\right)`$ (9) $`\xi _{\left(2\right)}^{ijkl}`$ $`=`$ $`{\displaystyle \frac{2\pi }{L^3}}{\displaystyle \underset{n}{}}\left(D_n^{(2)}\right)^{ij}\left(D_n^{(2)}\right)^{kl}\delta \left(2\omega \omega _n\right).`$ (10) The third term in (II), denoted $`\dot{n}_I`$, allows population control as discussed and observed by Fraser et al Fraser et al. (1999); Fraser and van Driel (2003). It can be written in terms of a third rank tensor $`\xi _{\left(I\right)}^{ijk}`$ as Fraser et al. (1999) $$\dot{n}_I=\xi _{\left(I\right)}^{ijk}E_\omega ^iE_\omega ^jE_{2\omega }^k+c.c.,$$ (11) where $$\xi _{(I)}^{ijk}=\frac{2\pi }{L^3}\underset{n}{}\left(D_n^{(2)}\right)^{ij}\left(D_n^{(1)}\right)^k\delta \left(2\omega \omega _n\right).$$ (12) It is purely real in the independent particle approximation Fraser et al. (1999); Fraser and van Driel (2003). Expressions such as (II) and (II) can also be written for carrier spin polarization and spin current Stevens et al. (2004). The interference terms of these describe ‘1+2’ spin control Stevens et al. (2003b) and ‘1+2’ spin current injection Bhat and Sipe (2000), which can be written in terms of material response pseudotensors $`\zeta _{(I)}^{ijkl}`$ Stevens et al. (2003b) and $`\mu _{(I)}^{ijklm}`$ Najmaie et al. (2003), respectively. In the independent particle approximation, $`\zeta _{(I)}^{ijkl}`$ is purely imaginary Stevens et al. (2003b), while $`\mu _{(I)}^{ijklm}`$ is purely real Najmaie et al. (2003). The phases of the material response tensors $`\eta _{(I)}`$, $`\xi _{(I)}`$, $`\mu _{(I)}`$, and $`\zeta _{(I)}`$ are related to the phases of the one- and two-photon matrix elements, $`𝐃_n^{(1)}`$ and $`𝖣_n^{\left(2\right)}`$. The one- and two-photon matrix elements also appear, respectively, in the one- and two-photon absorption coefficients, as can be seen from (9) and (10). There have been many theoretical investigations of one- and two-photon absorption near the direct gap of bulk semiconductors that include excitonic effects Dimmock (1967); Nathan et al. (1985). However, since one- and two-photon absorption are insensitive to the phases of the transition amplitudes, those calculations took no care to get the phases of the transition amplitudes correct. In the next two sections we find the transition amplitudes with the correct phases including excitonic effects. ## III Model We first review the two-band effective mass model of Wannier excitons; the two bands are nondegenerate conduction and valence bands that are parabolic and isotropic with a direct gap $`E_{cv}^g`$ at $`𝐤=\mathrm{𝟎}`$ (the $`\mathrm{\Gamma }`$ point) Yu and Cardona (1996); Bassani and Parravicini (1975). It has been used to study excitonic effects on one-photon absorption Elliott (1957), two-photon absorption Mahan (1968), and other nonlinear optical processes Blossey (1970); Ganguly and Birman (1967); Martin (1971); García-Cristóbal et al. (1994, 1998); Kolber and Dow (1978); Sheik-Bahae et al. (1994). We subsequently describe a generalization that accounts for degeneracy and multiple bands. It has been used for two-photon absorption Lee and Fan (1974), and has been implied whenever two-band results have been applied to actual semiconductors. The total Hamiltonian of the system can be written in the form $`H=H_B+H_C+H_{\text{int}}(t)`$. Here, $`H_0=H_B+H_C`$ is the field-free Hamiltonian made up of the single-particle part $`H_B`$ and the part due to the Coulomb interaction between carriers $`H_C`$. Using the minimal coupling Hamiltonian, the optical perturbation takes the form $`H_{\text{int}}(t)=\left(e/c\right)𝐀(t)𝐯+e^2A^2/(2mc^2)`$, where $`𝐀(t)`$ is the vector potential associated with the Maxwell electric field and $`𝐯`$ is the velocity operator associated with $`H_0`$. In the long wavelength limit, the position dependence of $`𝐀(t)`$ is neglected, and thus the second term in $`H_{\text{int}}(t)`$ may be neglected, since it can simply be absorbed in an overall time-dependent phase of the full system ket and hence cannot cause any transitions between states of $`H_0`$. Many approximate approaches to band structure calculation, including most pseudopotentials, and the truncation to a finite number of bands, implicitly assume an underlying field-free Hamiltonian that is nonlocal; there is then a correction to the interaction Hamiltonian $`H_{\text{int}}(t)`$ in the velocity gauge Girlanda et al. (1981); Aversa and Sipe (1995). However, we neglect such nonlocal corrections, as has been the practice in previous calculations of coherent control effects Atanasov et al. (1996); Fraser et al. (1999); Bhat and Sipe (2000). The initial state is the “vacuum” $`\|0`$; it corresponds to completely filled valence bands and empty conduction bands. If the Coulomb interaction were neglected in a two-band model consisting of valence ($`v`$) and conduction ($`c`$) bands, the final states would be of the form $`a_{c𝐤}^{}a_{v𝐤}\|0`$, where the operator $`a_{n𝐤}^{}`$ creates an electron in an eigenstate of $`H_B`$, a Bloch state $`\|n,𝐤`$ with band index $`n`$ and wavevector $`𝐤`$. The photon momentum has been neglected, consistent with the long wavelength approximation. The Coulomb interaction couples states at different $`𝐤`$; thus including $`H_C`$ the final states are of the form $$\|cv𝜿\underset{𝐤}{}A_{cv}^𝜿\left(𝐤\right)a_{c𝐤}^{}a_{v𝐤}\|0,$$ (13) where $`𝜿`$ labels the state; its physical meaning is given below. In the effective mass Wannier exciton approximation, the Fourier transform $$\psi _{cv}^𝜿(𝐫)=\underset{𝐤}{}A_{cv}^𝜿\left(𝐤\right)e^{i𝐤𝐫},$$ (14) which is the wavefunction for the relative coordinate between electron and hole, is a hydrogenic wavefunction satisfying $$\frac{\mathrm{}^2}{2\mu _{cv}}^2\psi _{cv}^𝜿(𝐫)V(r)\psi _{cv}^𝜿(𝐫)=\left(E_{cv}\left(𝜿\right)E_{cv}^g\right)\psi _{cv}^𝜿(𝐫),$$ (15) where $`\mu _{cv}^1=m_c^1+m_v^1`$ is the reduced mass in terms of the (positive) conduction and valence band effective masses, and $`V\left(r\right)`$ is the Coulomb potential, $`V(r)=e^2/\left(ϵr\right)`$, screened by the static dielectric constant $`ϵ`$ Elliott (1957); Baldereschi and Lipari (1971); Haug and Koch (1993). The state has energy $$E_{cv}\left(𝜿\right)=\frac{\mathrm{}^2\kappa ^2}{2\mu _{cv}}+E_{cv}^g.$$ We choose the states to be normalized over the volume $`L^3`$ by $`m,𝐤\|n,𝐤=\delta _{n,m}\delta _{𝐤,𝐤^{}}`$ and $`cv𝜿\|cv𝜿^{}=\delta _{𝜿,𝜿^{}}`$; as a result $`\psi _{cv}^𝜿(𝐫)`$ is unitless, having the normalization $`d^3r(\psi _{cv}^𝜿(𝐫))^{}\psi _{cv}^𝜿^{}(𝐫)=L^3\delta _{𝜿,𝜿^{}}`$. Our focus in this paper is on the unbound solutions to (15); bound exciton states lack relative velocity between the electron and hole, and hence do not contribute to the ballistic current or spin current. For a Fermi’s golden rule calculation of the current or spin current, the unbound state must behave asymptotically like an outgoing plane wave in the relative coordinate between electron and hole; $`𝜿`$ is the wavevector of the outgoing plane wave. Specifically, we must use “ionization states” rather than scattering states Breit and Bethe (1954), as was done for atomic ‘1+2’ ionization Baranova et al. (1990). They are related by $`\left(\psi _{cv}^𝜿(𝐫)\right)_{\text{ion}}=\left[\left(\psi _{cv}^𝜿(𝐫)\right)_{\text{scatt}}\right]^{}`$ Taylor (1972). Calculations of one- or two-photon absorption are insensitive to an error in this choice of boundary condition, but the present calculation is not, since it is sensitive to the relative phase of the transition amplitudes. The ionization state wavefunctions that solve (15) can be expressed as a sum over partial waves, $$\psi _{cv}^𝜿(𝐫)=e^{\frac{\pi }{2a_{cv}\kappa }}\underset{l=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }\left(l+1+\frac{i}{a_{cv}\kappa }\right)}{\left(2l\right)!}\left(2i\kappa r\right)^le^{i\kappa r}P_l\left(\frac{𝐫𝜿}{r\kappa }\right){}_{1}{}^{}F_{1}^{}(l+1+\frac{i}{a_{cv}\kappa };2l+2;2i\kappa r),$$ (16) where $`a_{cv}=ϵ\mathrm{}^2/\left(\mu _{cv}e^2\right)`$ is the exciton Bohr radius, and $`\kappa `$ and $`r`$ mean $`\left\|𝜿\right\|`$ and $`\left\|𝐫\right\|`$foo (a). The $`P_l`$ are Legendre polynomials, $`{}_{1}{}^{}F_{1}^{}`$ is a confluent hypergeometric function, and $`\mathrm{\Gamma }`$ is the Gamma function. Such a two-band model of Wannier excitons is useful for the description of many optical properties. However, near the band gap at the $`\mathrm{\Gamma }`$ point of a typical zincblende semiconductor there are, counting spin, eight bands: two each of conduction ($`c`$), heavy hole ($`hh`$), light hole ($`lh`$) and split-off hole ($`so`$). Other bands, especially the next higher conduction bands, can also be important for some processes, especially for population and spin control. The existence of multiple bands and band degeneracy modifies the exciton Hamiltonian, i.e., the operator acting on $`\psi _{cv}^𝜿(𝐫)`$ in the left side of (15). In the effective mass approximation, using a basis of $`\mathrm{\Gamma }`$ point states, the kinetic part of the Wannier exciton Hamiltonian has a matrix structure Luttinger and Kohn (1955); Dresselhaus (1956). Even though this approximation neglects band warping, nonparabolicity, and inversion asymmetry, the Hamiltonian lacks analytic eigenstates Dresselhaus (1956). This is essentially due to the degeneracy of the $`hh`$ and $`lh`$ bands at the $`\mathrm{\Gamma }`$ point. As a result of the difference between $`m_{hh}`$ and $`m_{lh}`$ there is ‘envelope-hole coupling,’ Sondergeld (1977b) which is a spin-orbit-like coupling between the orbital angular momentum of the exciton envelope function and the total angular momentum of the valence band $`\mathrm{\Gamma }`$ point Bloch functions Baldereschi and Lipari (1973). Baldereschi and Lipari split the effective mass Hamiltonian into a sum of terms based on symmetry, and showed that in a spherical approximation the envelope-hole coupling could be treated as a perturbation to the diagonal part, which has analytic, hydrogenic eigenstates Baldereschi and Lipari (1970, 1971). In order to extract the main physics, while preserving the simplicity of the two-band model, we neglect envelope-hole coupling entirely. In this approximation, (15) remains valid for each conduction-valence band pair, however one must use ‘average’ effective masses for degenerate bands. Specifically, the effective mass of the valence bands $`hh`$, $`lh`$, and $`so`$ is $`m/\gamma _{1L}`$, where $`m`$ is the free electron mass, and $`\gamma _{1L}`$ is one of the Luttinger parameters Baldereschi and Lipari (1971). The upper conduction bands have a different average effective mass. Note that $`\psi _{cv}^𝜿(𝐫)`$ is independent of $`c`$ and $`v`$ within the set of bands $`\{c,hh,lh,so\}`$. The effect of envelope-hole coupling has been studied for exciton bound states Baldereschi and Lipari (1970, 1971); Sondergeld (1977a), but not for optical processes involving unbound excitons in the continuum. Even within this model, the presence of multiple bands provides two types of terms in the sum over intermediate states in the two-photon amplitude: two-band terms, in which the intermediate and final states are in the same exciton series \[i.e., two states of the form (13) with the same $`c`$ and $`v`$\], and three-band terms, in which the intermediate and final states are in different series. Three-band terms are important for some processes but not for others. For current control, three-band terms are important for cross-linearly polarized fields Bhat and Sipe (2004), and for spin current control, they are important for the spin current due to co-linearly polarized fields Bhat and Sipe (2000). Three-band terms are essential for population and spin control Stevens et al. (2004); Bhat and Sipe (2004). The velocity matrix elements involving the state $`\|cv𝜿`$ are $$cv𝜿\left\|𝐯\right\|0=\underset{𝐤}{}\left(A_{cv}^𝜿\left(𝐤\right)\right)^{}𝐯_{cv}\left(𝐤\right)$$ (17) $$cv𝜿\left\|𝐯\right\|c^{}v^{}𝜿^{}=\underset{𝐤}{}\left(A_{cv}^𝜿\left(𝐤\right)\right)^{}A_{c^{}v^{}}^𝜿^{}\left(𝐤\right)\left(𝐯_{cc^{}}\left(𝐤\right)\delta _{v,v^{}}𝐯_{v^{}v}\left(𝐤\right)\delta _{c,c^{}}\right),$$ (18) where $`𝐯_{nm}\left(𝐤\right)n,𝐤\left\|𝐯\right\|m,𝐤`$ is the velocity matrix element between Bloch states. ## IV Transition amplitudes In the independent particle approximation, the transition amplitudes are $$\mathrm{\Omega }_{cv𝜿}^{(1\text{-free})}=i\frac{e}{2\mathrm{}\omega }𝐄_{2\omega }𝐯_{cv}\left(𝜿\right),$$ (19) and $`\mathrm{\Omega }_{cv𝜿}^{(2\text{-free})}=_{c^{},v^{}}\mathrm{\Omega }_{cc^{}vv^{}𝜿}^{(2\text{-free})}`$, where $$\mathrm{\Omega }_{cc^{}vv^{}𝜿}^{(2\text{-free})}\left(\frac{e}{\mathrm{}\omega }\right)^2\frac{\left(𝐄_\omega \left(𝐯_{cc^{}}\left(𝜿\right)\delta _{v,v^{}}\delta _{c,c^{}}𝐯_{v^{}v}\left(𝜿\right)\right)\right)\left(𝐄_\omega 𝐯_{c^{}v^{}}\left(𝜿\right)\right)}{E_{c^{}v^{}}/\mathrm{}\omega \left(𝜿\right)}.$$ (20) With excitonic effects included, using the perturbation $`H_{\text{int}}\left(t\right)`$ to second order gives the transition amplitudes $$\mathrm{\Omega }_{cv𝜿}^{(1)}=\frac{ie}{2\mathrm{}\omega }𝐄_{2\omega }cv𝜿\left\|𝐯\right\|0,$$ (21) and $$\mathrm{\Omega }_{cv𝜿}^{(2)}=\left(\frac{e}{\mathrm{}\omega }\right)^2\underset{c^{}v^{}𝜿^{}}{}\frac{\left(𝐄_\omega cv𝜿\left\|𝐯\right\|c^{}v^{}𝜿^{}\right)\left(𝐄_\omega c^{}v^{}𝜿^{}\left\|𝐯\right\|0\right)}{E_{c^{}v^{}}\left(𝜿^{}\right)/\mathrm{}\omega },$$ (22) where the sum over intermediate states is over both bound and free excitons. The two-photon transition amplitude is more difficult to deal with due to the sum over intermediate states; however, in our set of approximations it has been done exactly Mahan (1968); Lee and Fan (1974); Rustagi et al. (1973). In order to proceed analytically, it is common to use (17) and (18), and then make an expansion in $`𝐤`$ of the velocity matrix elements $`𝐯_{nm}\left(𝐤\right)`$ about the $`\mathrm{\Gamma }`$ point Elliott (1957); Loudon (1962); Mahan (1968); Rustagi et al. (1973); Lee and Fan (1974); Haug and Koch (1993). However, due to the degeneracy at the $`\mathrm{\Gamma }`$ point, the coefficients of such an expansion can depend on the direction of $`𝐤`$ Bir and Pikus (1974). To proceed, we note that Wannier excitons have large spatial extent and hence only a small region of wavevectors is important for them, i.e., $`A_{cv}^𝜿\left(𝐤\right)`$ is peaked in the region of $`𝐤`$ near $`𝜿`$ Bassani and Parravicini (1975). This is especially true for final states with energies above the band gap. Thus, we expand $`𝐯_{nm}\left(𝐤\right)`$ about the $`\mathrm{\Gamma }`$ point, approached in the direction $`\widehat{𝜿}`$, $$v_{nm}^i\left(𝐤\right)=v_{nm}^i\left(\widehat{𝜿}\right)+𝐤\mathbf{}_𝐤v_{nm}^i(\widehat{𝜿})+\mathrm{},$$ (23) where $`v_{nm}^i\left(\widehat{𝜿}\right)lim_{\lambda 0}n,\mathrm{\Gamma }\left\|𝐯\right\|m,\lambda 𝜿`$ and $`\mathbf{}_𝐤v_{nm}^i\left(\widehat{𝜿}\right)lim_{\lambda 0}\mathbf{}_𝜿n,\mathrm{\Gamma }\|𝐯m,\lambda 𝜿`$. Optical transitions due to the first term in (23) are called ‘allowed’, while those due to the second term are called ‘forbidden’. We restrict ourselves to materials for which the ‘allowed’ valence to conduction band transition does not vanish. Keeping only the ‘allowed’ term in (17),Elliott (1957) $$cv𝜿\left\|𝐯\right\|0=𝐯_{cv}\left(\widehat{𝜿}\right)\left(\psi _{cv}^𝜿(𝐫=\mathrm{𝟎})\right)^{}.$$ (24) For the intravalence and intraconduction band transitions, the first two terms of (23) in (18) give Rustagi et al. (1973) $$\begin{array}{cc}\hfill cv𝜿\left\|𝐯\right\|c^{}v^{}𝜿^{}=& \left[\delta _{v,v^{}}𝐯_{cc^{}}\left(\widehat{𝜿}\right)\delta _{c,c^{}}𝐯_{v^{}v}\left(\widehat{𝜿}\right)\right]\frac{d^3r}{L^3}\left(\psi _{cv}^𝜿\left(𝐫\right)\right)^{}\psi _{c^{}v^{}}^𝜿^{}(𝐫)\hfill \\ & \left[\delta _{v,v^{}}\mathbf{}_𝐤v_{cc^{}}^i\left(\widehat{𝜿}\right)\delta _{c,c^{}}\mathbf{}_𝐤v_{v^{}v}^i\left(\widehat{𝜿}\right)\right]i\frac{d^3r}{L^3}\left(\psi _{cv}^𝜿\left(𝐫\right)\right)^{}\mathbf{}\psi _{c^{}v^{}}^𝜿^{}(𝐫).\hfill \end{array}$$ (25) In particular Rustagi et al. (1973); Bethe and Salpeter (1977), $$cv𝜿\left\|𝐯\right\|cv𝜿=\underset{𝐤}{}\left(A_{cv}^\kappa \left(𝐤\right)\right)^{}A_{cv}^\kappa \left(𝐤\right)\frac{\mathrm{}}{\mu _{cv}}𝐤=i\frac{\mathrm{}}{\mu _{cv}}\frac{d^3r}{L^3}\left(\psi _{cv}^\kappa (𝐫)\right)^{}_𝐫\psi _{cv}^\kappa (𝐫)=\mathrm{}𝜿/\mu _{cv}.$$ (26) For Ge and for simple models of zincblende semiconductors that neglect lack of inversion symmetry, the first term in (25) always vanishes. This means that there are only ‘allowed-forbidden’ two-photon transitions (the interband transition is allowed, while the intraband transition is forbidden). When the first term is nonvanishing, there are ‘allowed-allowed’ two-photon transitions. In principle, for materials that lack a center of inversion, there is also a small contribution to the ‘allowed-forbidden’ two-photon transition from the first term in (25) and the term in $`cv𝜿\left\|𝐯\right\|0`$ that comes from the second term in (23); we neglect it in what follows, but note that when compared to the dominant ‘allowed-forbidden’ contribution that we consider here, it has a different Coulomb enhancement but the same intrinsic phase (see Eq. 2.32 of Rustagi et alRustagi et al. (1973)). We write $`\mathrm{\Omega }_{cv𝜿}^{(2)}=\mathrm{\Omega }_{cv𝜿}^{(2\text{:a-f})}+\mathrm{\Omega }_{cv𝜿}^{(2\text{:a-a})}`$, and discuss the ‘allowed-forbidden’ and ‘allowed-allowed’ transitions separately. Using (16), (24), and (21), the one-photon transition amplitude is Elliott (1957) $`\mathrm{\Omega }_{cv𝜿}^{(1)}`$ $`=`$ $`\mathrm{\Omega }_{cv𝜿}^{(1\text{-free})}\mathrm{exp}\left({\displaystyle \frac{\pi }{2a_{cv}\kappa }}\right)\mathrm{\Gamma }\left(1{\displaystyle \frac{i}{a_{cv}\kappa }}\right),`$ (27) where only the ‘allowed’ term is kept in $`\mathrm{\Omega }_{cv𝜿}^{(1\text{-free})}`$. The transition is to an $`s`$-wave. The one-photon absorption coefficient is proportional to the norm of $`\mathrm{\Omega }_{cv𝜿}^{(1)}`$ \[see (II) and (9)\])Elliott (1957). For ‘allowed-forbidden’ two-photon transitions, substituting (24) and the second term of (25) into (22), $$\mathrm{\Omega }_{cv𝜿}^{(2\text{:a-f})}=\frac{e^2}{\omega ^2\mathrm{}}\underset{c^{}v^{}}{}\left(𝐄_\omega 𝐯_{c^{}v^{}}\left(\widehat{𝜿}\right)\right)E_\omega ^i\left(\delta _{v,v^{}}\mathbf{}_𝐤v_{cc^{}}^i\left(\widehat{𝜿}\right)\delta _{c,c^{}}\mathbf{}_𝐤v_{v^{}v}^i\left(\widehat{𝜿}\right)\right)𝐌_{cc^{}vv^{}}\left(𝜿\right),$$ (28) where $$𝐌_{cc^{}vv^{}}\left(𝜿\right)id^3r\left(\psi _{cv}^𝜿\left(𝐫\right)\right)^{}\mathbf{}G_{c^{}v^{}}(𝐫,\mathrm{𝟎};\mathrm{}\omega E_{c^{}v^{}}^g),$$ (29) and we have used the Coulomb Green function, $$G_{cv}(𝐫,𝐫^{};\mathrm{\Omega })=\frac{1}{L^3}\underset{𝜿}{}\frac{\psi _{cv}^𝜿(𝐫)\left(\psi _{cv}^𝜿(𝐫^{})\right)^{}}{E_{cv}\left(\kappa \right)E_{cv}^g\mathrm{\Omega }},$$ which is known analytically Mahan (1968). In particular, $$G_{cv}(𝐫,\mathrm{𝟎};\mathrm{}\omega E_{cv}^g)=\frac{\mu _{cv}}{2\pi r\mathrm{}^2}\mathrm{\Gamma }\left(1\gamma _{cv}\right)W_{\gamma _{cv},\frac{1}{2}}\left(\frac{2r}{a_{cv}\gamma _{cv}}\right),$$ (30) where we define $$\gamma _{cv}\sqrt{\frac{B_{cv}}{E_{cv}^g\mathrm{}\omega }},$$ $`B_{cv}=\mathrm{}^2/\left(2\mu _{cv}a_{cv}^2\right)`$ is the exciton binding energy, and $`W_{\gamma ,1/2}\left(z\right)`$ is a Whittaker function with the integral representation $$W_{\gamma ,1/2}\left(z\right)=\frac{ze^{\frac{z}{2}}}{\mathrm{\Gamma }\left(1\gamma \right)}_0^{\mathrm{}}𝑑t\left(\frac{1+t}{t}\right)^\gamma e^{zt}.$$ (31) Since the Green function depends only on the magnitude of $`𝐫`$, only the p-wave of the final state survives the integral in Eq. (29) over the angles of $`𝐫`$, $`𝑑\mathrm{\Omega }P_1\left(\widehat{𝐫}\widehat{𝜿}\right)\widehat{𝐫}=4\pi \widehat{𝜿}/3`$. The integral over $`r`$ can be done using Mahan (1968) $$_0^{\mathrm{}}r^{\sigma 1}e^{pr}{}_{1}{}^{}F_{1}^{}(\alpha ;\sigma ;\lambda r)𝑑r=\mathrm{\Gamma }\left(\sigma \right)\frac{p^{\alpha \sigma }}{\left(p\lambda \right)^\alpha }.$$ (32) The final result is $$\mathrm{\Omega }_{cv𝜿}^{(2\text{:a-f})}=\mathrm{exp}\left(\frac{\pi }{2a_{cv}\kappa }\right)\mathrm{\Gamma }\left(2\frac{i}{a_{cv}\kappa }\right)\underset{c^{}v^{}}{}N_{cc^{}vv^{}}^{(\text{a-f})}\left(\kappa \right)\mathrm{\Omega }_{cc^{}vv^{}𝜿}^{(2\text{-free})},$$ (33) where only the ‘allowed-forbidden’ term is kept in $`\mathrm{\Omega }_{cc^{}vv^{}𝜿}^{(2\text{-free})}`$, and $$N_{cc^{}vv^{}}^{(\text{a-f})}\left(\kappa \right)\left(1+a_{c^{}v^{}}^2\kappa ^2\gamma _{c^{}v^{}}^2\right)2_0^1\frac{S\left(\frac{1+S}{1S}\right)^{\gamma _{c^{}v^{}}}\mathrm{exp}\left(\frac{2}{a_{cv}\kappa }\mathrm{arctan}\left(a_{c^{}v^{}}\kappa \gamma _{c^{}v^{}}S\right)\right)}{\left(1+a_{c^{}v^{}}^2\kappa ^2\gamma _{c^{}v^{}}^2S^2\right)^2}𝑑S.$$ (34) For ‘allowed-allowed’ two-photon transitions, substituting (24) and the first term of (25) into (22), $$\begin{array}{cc}\hfill \mathrm{\Omega }_{cv𝜿}^{(2\text{:a-a})}=& \left(\frac{e}{\mathrm{}\omega }\right)^2\underset{c^{}v^{}}{}𝐄_\omega \left[\delta _{vv^{}}v_{cc^{}}^i\left(\widehat{𝜿}\right)\delta _{c,c^{}}v_{v^{}v}^i\left(\widehat{𝜿}\right)\right]𝐄_\omega 𝐯_{c^{}v^{}}\left(\widehat{𝜿}\right)\hfill \\ & \times \mathrm{}d^3r\left(\psi _{cv}^𝜿\left(𝐫\right)\right)^{}G_{c^{}v^{}}(𝐫,\mathrm{𝟎};\mathrm{}\omega E_{c^{}v^{}}^g).\hfill \end{array}$$ Since $`G_{c^{}v^{}}`$ depends only on the magnitude of $`𝐫`$ \[Eq. (30)\], only the $`s`$ part of the final state will survive the integration over angles of $`𝐫`$. Again we use (31) for the Whittaker function. The integral over the magnitude of $`𝐫`$ can be done using an identity obtained by taking a derivative with respect to $`p`$ of both sides of (32). Finally, $$\mathrm{\Omega }_{cv𝜿}^{(2\text{:a-a})}=\mathrm{exp}\left(\frac{\pi }{2a_{cv}\kappa }\right)\mathrm{\Gamma }\left(1\frac{i}{a_{cv}\kappa }\right)\underset{c^{}v^{}}{}\mathrm{\Omega }_{cc^{}vv^{}𝜿}^{(2\text{-free})}N_{cc^{}vv^{}}^{(\text{a-a})}\left(\kappa \right),$$ (35) where only the ‘allowed-allowed’ term is kept in $`\mathrm{\Omega }_{cc^{}vv^{}𝜿}^{(2\text{-free})}`$, and $$\begin{array}{cc}\hfill N_{cc^{}vv^{}}^{(\text{a-a})}\left(\kappa \right)& \left(1+\left(a_{c^{}v^{}}\kappa \gamma _{c^{}v^{}}\right)^2\right)2\hfill \\ & \times _0^1S(1S\frac{a_{c^{}v^{}}\gamma _{c^{}v^{}}}{a_{cv}})\left(\frac{1+S}{1S}\right)^{\gamma _{c^{}v^{}}}\frac{\mathrm{exp}\left(\frac{2}{a_{cv}\kappa }\mathrm{arctan}\left(a_{c^{}v^{}}\kappa \gamma _{c^{}v^{}}S\right)\right)}{\left(1+\left(a_{c^{}v^{}}\kappa \gamma _{c^{}v^{}}S\right)^2\right)^2}dS.\hfill \end{array}$$ (36) This agrees with Eq. 2.28 of Rustagi Rustagi et al. (1973), but note that we have defined $`N_{cc^{}vv^{}}^{(\text{a-a})}\left(\kappa \right)=\left(1+\left(a_{c^{}v^{}}\kappa \gamma _{c^{}v^{}}\right)^2\right)I_{s,k}\left(\kappa \right)`$, where $`I_{s,k}\left(\kappa \right)`$ is given, with a typographical error, in Eq. 2.25 of that paper. The factors $`N_{cc^{}vv^{}}^{(\text{a-f})}`$ and $`N_{cc^{}vv^{}}^{(\text{a-a})}`$, which appear in (34) and (36) are the enhancements due to the Coulomb interaction in the intermediate states; they are discussed further in Appendix A. ## V Results The one- and two-photon transition amplitudes were presented in the previous section on the basis of an expansion in $`𝐤`$ of the Bloch state velocity matrix elements. The ‘allowed’ one-photon transition amplitude $`\mathrm{\Omega }^{(1)}`$ is in (27), the ‘allowed-forbidden’ two-photon transition amplitude is in (33), and the ‘allowed-allowed’ two-photon transition amplitude is in (35). From them, $`𝐃_{cv𝜿}^{\left(1\right)}`$ and $`𝖣_{cv𝜿}^{\left(2\right)}`$ may be extracted by comparison with the definitions in (5) and (6). ### V.1 Current injection The ‘1+2’ current injection is dominated by interference of ‘allowed’ one-photon transitions and ‘allowed-forbidden’ two-photon transitions Bhat and Sipe (2004). Substituting these into (7), and using the Gamma function identities $`\mathrm{\Gamma }\left(x+1\right)=x\mathrm{\Gamma }\left(x\right)`$ and $$\mathrm{\Gamma }\left(1ix\right)\mathrm{\Gamma }\left(1+ix\right)=\pi x/\mathrm{sinh}\left(\pi x\right),$$ (37) yields our final result for the current injection tensor $$\eta _{(I)}^{ijkl}=\underset{c,v}{}\left(1+\frac{i}{a_{cv}\kappa _{cv}}\right)\mathrm{\Xi }\left(a_{cv}\kappa _{cv}\right)\underset{c^{},v^{}}{}N_{cc^{}vv^{}}^{(\text{a-f})}\left(\kappa _{cv}\right)\eta _{cc^{}vv^{}}^{ijkl},$$ (38) where $$\kappa _{cv}\frac{1}{a_{cv}}\sqrt{\frac{2\mathrm{}\omega E_{cv}^g}{B_{cv}}},$$ (39) $$\mathrm{\Xi }\left(x\right)\frac{\left(\pi /x\right)\mathrm{exp}\left(\pi /x\right)}{\mathrm{sinh}\left(\pi /x\right)}=\frac{2\pi }{x}\left(1\mathrm{exp}\left(2\pi /x\right)\right)^1,$$ (40) and $$\eta _{cc^{}vv^{}}^{ijkl}\frac{2\pi e}{L^3}\underset{𝜿}{}\mathrm{\Delta }_{cv}^i\left(𝜿\right)\left(D_{cc^{}vv^{}𝜿}^{(2\text{-free})}\right)^{jk}\left(D_{cv𝜿}^{(1\text{-free})}\right)^l\delta \left(2\omega E_{cv}\left(\kappa \right)/\mathrm{}\right),$$ (41) with $`𝚫_{cv}\left(𝜿\right)\left(𝐯_{cc}\left(𝜿\right)𝐯_{vv}\left(𝜿\right)\right)=\mathrm{}𝜿/\mu _{cv}`$ and $$\left(D_{cc^{}vv^{}𝜿}^{(2\text{-free})}\right)^{jk}=\left(\frac{e}{\mathrm{}\omega }\right)^2\frac{\{\left(𝐯_{cc^{}}\left(𝜿\right)\delta _{v,v^{}}\delta _{c,c^{}}𝐯_{v^{}v}\left(𝜿\right)\right),𝐯_{c^{}v^{}}\left(𝜿\right)\}^{jk}}{E_{c^{}v^{}}\left(\kappa \right)/\mathrm{}\omega },$$ (42) where $`\{𝐯_1,𝐯_2\}^{ij}(v_1^iv_2^j+v_1^jv_2^i)/2`$ and $$\left(D_{cv𝜿}^{(1\text{-free})}\right)^l=i\frac{e}{2\mathrm{}\omega }v_{cv}^l\left(𝜿\right).$$ (43) Note that only the ‘allowed’ part of (43) and the ‘allowed-forbidden’ part of (42) should be retained for a consistent solution. We have written (38) to separate the parts due to the electron-hole interaction. In the independent particle approximation, the current injection tensor $`\eta _{(I\text{-free})}^{ijkl}`$ is Atanasov et al. (1996) $$\eta _{(I\text{-free})}^{ijkl}=\underset{c,c^{},v,v^{}}{}\eta _{cc^{}vv^{}}^{ijkl};$$ (44) it is evaluated for parabolic bands in Appendix B. For GaAs, we present in Fig. 1 the magnitude of $`\eta _{(I)}^{xxxx}`$, based on $`\eta _{cc^{}vv^{}}^{xxxx}`$ calculated by two methods. The first method, described in Appendix B, uses isotropic parabolic bands and includes only two-band terms; it uses effective mass ratios for conduction, heavy hole, light hole, and split-off bands of $`0.067`$, $`0.51`$, $`0.082`$, and $`0.154`$ respectively, $`E_P=27.86`$ eV, the fundamental band gap $`E_g`$ is 1.519 eV, and valence band spin-orbit splitting is 0.341 eV Madelung (1996); Pfeffer and Zawadzki (1990). The second method solves the $`8\times 8`$ $`𝐤𝐩`$ Hamiltonian including remote band effects, but in a spherical approximation with warping and spin-splitting neglected by replacing $`\gamma _2`$ and $`\gamma _3`$ with $`\stackrel{~}{\gamma }\left(2\gamma _2+3\gamma _3\right)/5`$ Baldereschi and Lipari (1973); the calculation is nonperturbative in $`𝐤`$ (hence it includes band nonparabolicity) and it includes both two- and three-band terms in the two-photon amplitude Bhat and Sipe (2004). The solid and dotted lines in Fig. 1 are calculated with (38), and hence include excitonic effects; the Coulomb enhancement part of the calculation uses $`B_{cv}=4.2`$ meV Sell (1972) and the band parameters listed above. Note that the solid black line in Fig. 1 is inconsistent in the sense that the Coulomb enhancement is based on an expansion in $`𝐤`$, whereas the free-particle result that it enhances is nonperturbative in $`𝐤`$; nevertheless, such an approach has given good agreement with experiments for one- and two-photon absorption Sturge (1962); Weiler (1981). The Coulomb enhancement of $`\eta _{(I)}`$ can clearly be seen in Fig. 1. There is a kink in each curve at excess photon energy 341 meV corresponding to the onset of transitions from the $`so`$ band. At higher energies, the Coulomb enhancement of $`so`$ transitions is larger than the Coulomb enhancements of $`hh`$ and $`lh`$ transitions, since the former transitions are to conduction band states with lower energy. Hence, the kink in $`\eta _{(I)}`$ is enhanced by excitonic effects. We extract the intrinsic phase of $`\eta _{(I)}^{xxxx}`$ using (2). The solid line in Fig. 2 is the intrinsic phase of $`\eta _{(I)}^{xxxx}`$ calculated for GaAs with the nonperturbative $`8\times 8`$ $`𝐤𝐩`$ Hamiltonian band model; the result for the parabolic band model is nearly identical. Since we have used a spherical exciton model, the intrinsic phase is the same for all components of $`\eta _{(I)}^{ijkl}`$. The intrinsic phase has its maximum value of $`\pi /2`$ at the band edge, and goes to zero as the light frequency increases. The decrease is smooth except for a small kink at the onset of transitions from the $`so`$ band. In fact, for excess photon energies less than the split-off energy, the intrinsic phase has the simple analytic form $$\delta \left(\omega \right)=\mathrm{arctan}\left(\sqrt{\frac{B_{cv}}{2\mathrm{}\omega E_g}}\right).$$ (45) Equation (45) is plotted as the dotted line in Fig. 2; compared to the solid line, it is identical below the onset of $`so`$ transitions, and it makes a good approximation above the the onset of $`so`$ transitions. Since (45) only depends on the excess photon energy scaled by the exciton binding energy, we plot it as a function of this scaled energy in the inset of Fig. 2; it is useful for finding the intrinsic phase of materials other than GaAs. In $`\eta _{(I)}`$ \[Eq. (38)\], the two- and three-band terms have different intermediate state Coulomb enhancement $`N_{cc^{}vv^{}}^{(\text{a-f})}`$. For many materials, however, $`N_{cc^{}vv^{}}^{(\text{a-f})}`$ is approximately equal for all the terms $`\eta _{cc^{}vv^{}}^{ijkl}`$ that contribute significantly to the total $`\eta _{(I\text{-free})}^{ijkl}`$, as shown in Appendix A for GaAs. Thus, at photon energies for which transitions from the heavy- and light-hole bands dominate $`\eta `$, the Coulomb enhancement becomes approximately independent of the sum over bands and we can make the simplification $$\eta _{(I)}^{ijkl}F_{\text{a-f}}^{(I)}\mathrm{exp}\left(i\delta \right)\eta _{(I\text{-free})}^{ijkl},$$ (46) where the intrinsic phase is given by (45), and $$F_{\text{a-f}}^{(I)}\left(\omega \right)\mathrm{\Xi }\left(a_{cv}\kappa _{cv}\right)\sqrt{1+\left(a_{cv}\kappa _{cv}\right)^2}N_{ccvv}^{(\text{a-f})}\left(\kappa _{cv}\right),$$ (47) The Coulomb enhancement factor $`F_{\text{a-f}}^{(I)}\left(\omega \right)`$ is plotted in Fig. 3 with the approximation that $`N_{ccvv}^{(\text{a-f})}=1`$ (see Appendix A). ### V.2 Carrier population control The ‘1+2’ carrier population control is dominated by interference of ‘allowed’ one-photon transitions and ‘allowed-allowed’ two-photon transitions Bhat and Sipe (2004); Fraser and van Driel (2003); Stevens et al. (2004). Substituting these into (12), and using the Gamma function identity (37), we find $$\xi _{\left(I\right)}^{ijk}=\underset{c,v}{}\mathrm{\Xi }\left(a_{cv}\kappa _{cv}\right)\underset{c^{}v^{}}{}N_{cc^{}vv^{}}^{(\text{a-a})}\left(\kappa _{cv}\right)\xi _{cc^{}vv^{}}^{ijk},$$ (48) where $$\xi _{cc^{}vv^{}}^{ijk}=\frac{2\pi e}{L^3}\underset{𝜿}{}\left(D_{cc^{}vv^{}𝜿}^{(2\text{-free})}\right)^{jk}\left(D_{cv𝜿}^{(1\text{-free})}\right)^l\delta \left(2\omega E_{cv}\left(\kappa \right)/\mathrm{}\right)$$ (49) and $`𝖣_{cc^{}vv^{}𝜿}^{(2\text{-free})}`$ and $`𝐃_{cv𝜿}^{(1\text{-free})}`$ are given by (42) and (43). Note that only the ‘allowed’ part of (43) and the ‘allowed-allowed’ part of (42) should be retained for a consistent solution. In the independent particle approximation, $$\xi _{\left(I\text{-free}\right)}^{ijk}=\underset{c^{}v^{}}{}\xi _{cc^{}vv^{}}^{ijk}.$$ (50) Thus, population control has a Coulomb enhancement due to excitonic effects, but no phase shift. Note that (48) gives the population control tensor at final energies above the band edge. There can also be population control of bound excitons when both one- and two-photon transitions are to the same excitonic state. This can occur, for example, at $`s`$ excitons due to allowed-allowed two-photon transitions Doni et al. (1974) interfering with allowed one-photon transitions. If $`N_{cc^{}vv^{}}^{(\text{a-a})}\left(\kappa _{cv}\right)`$ is approximately the same for all the terms that significantly contribute to $`\xi _{\left(I\right)}`$, then, at photon energies for which transitions from the heavy and light hole bands dominate $`\xi _{(I)}`$, the Coulomb enhancement becomes approximately independent of the sum over bands, and we can make the simplification, $$\xi _{\left(I\right)}^{ijk}F_{\text{a-a}}^{(I)}\xi _{\left(I\text{-free}\right)}^{ijk},$$ (51) where $$F_{\text{a-a}}^{(I)}\left(\omega \right)\mathrm{\Xi }\left(a_{cv}\kappa _{cv}\right)N_{ccvv}^{(\text{a-a})}\left(\kappa _{cv}\right).$$ (52) The Coulomb enhancement factor $`F_{\text{a-a}}^{(I)}\left(\omega \right)`$ is plotted in Fig. 3 with the approximation that $`N_{ccvv}^{(\text{a-a})}=1`$ (see Appendix A). ### V.3 Spin current injection and spin control The ‘1+2’ spin current is dominated by interference of ‘allowed’ one-photon transitions and ‘allowed-forbidden’ two-photon transitions, whereas ‘1+2’ spin control is dominated by interference of ‘allowed’ one-photon transitions and ‘allowed-allowed’ two-photon transitions Bhat and Sipe (2004). Under the approximations that led to (46) and (51), the spin current injection pseudotensor is $$\mu _{(I)}^{ijklm}=F_{\text{a-f}}^{(I)}\mathrm{exp}\left(i\delta \right)\mu _{(I\text{-free})}^{ijklm},$$ (53) where $`F_{\text{a-f}}^{(I)}`$ is given by (47), $`\delta `$ is given by (45), and $`\mu _{(I\text{-free})}`$ is the spin current injection pseudotensor in the independent particle approximation. Under similar approximations, the spin control pseudotensor is $$\zeta _{(I)}^{ijkl}=F_{\text{a-a}}^{(I)}\zeta _{(I\text{-free})}^{ijkl},$$ (54) where $`F_{\text{a-a}}^{(I)}`$ is given by (52), and $`\zeta _{(I\text{-free})}`$ is the spin control pseudotensor in the independent particle approximation. Spin control, like carrier population control, has a Coulomb enhancement but no phase shift. There can also be spin control of bound excitons, but it has not been included in (54). ## VI Discussion We now examine the relationship between the Coulomb enhancements of the ‘1+2’ processes and of one- and two-photon absorption; the latter are denoted by $`F^{\left(1\right)}`$ and $`F^{\left(2\right)}`$ so that for $`i\{1,2\}`$, $`\dot{n}^{\left(i\right)}=\dot{n}_{\text{free}}^{\left(i\right)}F^{\left(i\right)}`$. The relationship is particularly simple at photon energies for which transitions from the heavy- and light-hole bands are dominant and intermediate state Coulomb enhancement is the same for each significant term in the sum over intermediate states. The Coulomb enhancements for the ‘1+2’ processes are then given by (47) and (52). For one-photon absorption, $`F^{\left(1\right)}=\mathrm{\Xi }\left(a_{cv}\kappa _{cv}\right)`$ Elliott (1957). In noncentrosymmetric semiconductors, two-photon absorption is dominated by allowed-allowed transitions just above the band gap, and by allowed-forbidden transitions at higher final energies; the cross-over point in GaAs is a few meV above the band gap van der Ziel (1977). At photon energies for which allowed-allowed transitions dominate two-photon absorption, from (35), $$F^{\left(2\right)}=\mathrm{\Xi }\left(a_{cv}\kappa _{cv}\right)\left(N_{ccvv}^{(\text{a-a})}\left(\kappa _{cv}\right)\right)^2,$$ (55) and thus $$F_{\text{a-a}}^{(I)}=\sqrt{F^{\left(1\right)}F^{\left(2\right)}}\text{ and }F_{\text{a-f}}^{(I)}=C\sqrt{F^{\left(1\right)}F^{\left(2\right)}},$$ (56) where $`C\left(N_{ccvv}^{(\text{a-f})}\left(\kappa _{cv}\right)/N_{ccvv}^{(\text{a-a})}\left(\kappa _{cv}\right)\right)\sqrt{1+\left(a_{cv}\kappa _{cv}\right)^2}`$, while at photon energies for which allowed-forbidden transitions dominate two-photon absorption, from (33), $$F^{\left(2\right)}=\mathrm{\Xi }\left(a_{cv}\kappa _{cv}\right)\left(1+\left(a_{cv}\kappa _{cv}\right)^2\right)\left(N_{ccvv}^{(\text{a-f})}\left(\kappa _{cv}\right)\right)^2,$$ (57) and thus $$F_{\text{a-a}}^{(I)}=\left(1/C\right)\sqrt{F^{\left(1\right)}F^{\left(2\right)}}\text{ and }F_{\text{a-f}}^{(I)}=\sqrt{F^{\left(1\right)}F^{\left(2\right)}}.$$ (58) Note that, based on Appendix A, $`C\sqrt{1+B_{cv}/\left(2\mathrm{}\omega E_g\right)}`$, which is the ratio of the two curves in Fig. 3. In centrosymmetric semiconductors, there are no allowed-allowed transitions, and only (58) applies. The ‘1+2’ processes are often described by ratios. For example, a useful quantity to describe the current is the swarm velocity Haché et al. (1998); Sipe et al. (1998), defined as the average velocity per injected electron-hole pair $$𝐯_{\text{swarm}}\frac{\left(d𝐉/dt\right)}{e\left(dn/dt\right)}.$$ The swarm velocity is a maximum when the relative intensities of the two colors are chosen such that $`\dot{n}_{2\omega }=\dot{n}_\omega `$; returning to (II), if one associates the one- and two-photon amplitudes with the arms of an effective interferometer, this condition corresponds to balancing that interferometer. For fields co-linearly polarized along $`\widehat{𝐱}`$, the maximum swarm speed is $$v_{\text{swarm}}=\frac{1}{e}\frac{\left\|\eta _{(I)}^{xxxx}\right\|}{\sqrt{\xi _{\left(1\right)}^{xx}\xi _{\left(2\right)}^{xxxx}}}.$$ (59) A useful quantity to describe pure spin currents is the maximum spin separation distance Hübner et al. (2003); it is proportional to $`\mu _{(I)}/\sqrt{\xi _{(1)}\xi _{(2)}}`$. As a consequence of (58), the maximum swarm speed, and the maximum spin separation distance, are unaffected by excitonic effects when allowed-forbidden transitions dominate two-photon absorption foo (b). However, close to the band edge, where allowed-allowed transitions dominate two-photon absorption, excitonic effects increase these ratios by a factor $`C`$ over their value in the independent particle approximation. In contrast, as a consequence of (58), excitonic effects do not affect the maximum control ratio for population and spin control ($`\xi _{(I)}/\sqrt{\xi _{(1)}\xi _{(2)}}`$ and $`\zeta _{(I)}/\sqrt{\xi _{(1)}\xi _{(2)}}`$ respectively Fraser et al. (1999); Stevens et al. (2003b); Stevens et al. (2004)) close to the band edge and decrease them by a factor $`C`$ at higher photon energies for which allowed-forbidden transitions dominate two-photon absorption. In the terminology of Seideman Seideman (1998), the excitonic phase shift of the ‘1+2’ current and spin current is a direct phase shift. It is due to the complex nature of the final state as it appears in the transition amplitudes. Thus it can be understood in terms of the partial wave phase shifts of the final state caused by the Coulomb potential between electron and hole. The Coulomb interaction is rather unique due to its long range nature, so we first suppose the potential between the electron and hole falls off more rapidly than $`1/C`$. In that simpler problem, the final state wave function is written as $$\psi _𝜿(𝐫)=\underset{l=0}{\overset{\mathrm{}}{}}i^le^{i\delta _l\left(\kappa \right)}\left(2l+1\right)\frac{u_{\kappa ,l}\left(r\right)}{r}P_l\left(\frac{𝐫𝜿}{r\kappa }\right),$$ where the $`u_{\kappa ,l}\left(r\right)`$ are real Taylor (1972). If the potential between the particles is ignored, then the partial wave phase shifts, $`\delta _l\left(\kappa \right)`$ are zero. The allowed one-photon pathway reaches an $`s`$ wave, while the allowed-forbidden two-photon pathway reaches a $`p`$ wave. Substituting this form for the wave function into the one- and two-photon transition amplitudes, yields $`\mathrm{\Omega }_𝜿^{(1)}=\mathrm{\Omega }_𝜿^{(1\text{-free})}e^{i\delta _0\left(\kappa \right)}f_0\left(\kappa \right)`$ for the one-photon rate, where $`f_0\left(\kappa \right)`$ is real and depends on $`u_{\kappa ,0}\left(r\right)`$, and $`\mathrm{\Omega }_𝜿^{(2\text{:a-f})}=\mathrm{\Omega }_𝜿^{(2\text{-free})}e^{i\delta _1\left(\kappa \right)}f_1\left(\kappa \right)`$ for the two-photon rate, where $`f_1\left(\kappa \right)`$ is real and depends on $`u_{\kappa ,0}\left(r\right)`$ and $`u_{\kappa ,1}\left(r\right)`$. Here $`\mathrm{\Omega }_𝜿^{(i\text{-free})}`$ is the i-photon transition amplitude when the potential between the particles is ignored. It is then straightforward to see from (7) that the relative shift of the partial waves is responsible for the phase shift of the current and spin current. That is, $$\delta =\delta _0\delta _1.$$ (60) The use of ionization states as opposed to scattering states was important to get the correct sign of the intrinsic phase. With scattering states, one would find $`\delta =\delta _1\delta _0`$. Due to the long-range nature of the Coulomb potential, the partial wave phase shifts have a logarithmic $`r`$ dependent part, but it is the same for all partial waves and thus does not appear in the relative phase. The part of the Coulomb partial wave phase shift $`\delta _l\left(\kappa \right)`$ that does not depend on $`r`$ is $`\mathrm{arg}\left(\mathrm{\Gamma }\left(l+1+i/(a_{cv}\kappa )\right)\right)`$ Bethe and Salpeter (1977); Taylor (1972); when inserted into (60), this reproduces (45). In contrast, the allowed-allowed two-photon pathway reaches an $`s`$ wave and thus there is no phase shift for population control or spin control. The expression (60) for the intrinsic phase in terms of the scattering phases is particularly simple, since each pathway connects to only a single parity. This contrasts with ‘1+2’ ionization from an atomic $`s`$ state, for which the one-photon transition is to a $`p`$ wave and the two-photon transition is to both $`s`$ and $`d`$ waves; the intrinsic phase is thus a weighting of the $`p`$-$`s`$ and $`p`$-$`d`$ partial wave shifts Baranova et al. (1990). Materials for which the first term in (23) is forbidden (Cu<sub>2</sub>O is an example) have these same selection rules Elliott (1957); Mahan (1968); Rustagi et al. (1973); hence, they will have an intrinsic phase with a similar weighting. The absence of a phase shift in population control can be connected to a symmetry of the second order nonlinear optical susceptibility. From considerations of energy transfer and macroscopic electrodynamics, $`\xi _{(I)}`$ is related to the nonlinear susceptibility $`\chi ^{(2)}`$ by $$\xi _{(I)}^{ijk}=\left(iϵ_0/\mathrm{}\right)\left[\chi ^{(2)kij}(2\omega ;\omega ,\omega )\chi ^{(2)jki}(\omega ;2\omega ,\omega )\right].$$ (61) In the independent particle approximation Sipe and Shkrebtii (2000), $$\chi ^{\left(2\right)ijk}(2\omega ;\omega ,\omega )=\left[\chi ^{\left(2\right)jik}(\omega ;2\omega ,\omega )\right]^{},$$ (62) which is a generalization of overall permutation symmetry to resonant absorption. As a result of (62), Fraser et al. showed that $`\xi _{\left(I\right)}`$ is proportional to $`\mathrm{Im}\chi ^{(2)}`$, and is thus purely real Fraser et al. (1999). Our result that $`\xi _{(I)}`$ remains real when excitonic effects are included suggests that (62) holds more generally. In fact, it can be shown that (62) holds for any Hamiltonian symmetric under time-reversal so long as $`\mathrm{}\omega `$ is not resonant. ## VII Summary and Outlook We have extended the theory of interband ‘1+2’ processes in bulk semiconductors to include the electron-hole interaction. Following previous theories Atanasov et al. (1996); Fraser et al. (1999); Bhat and Sipe (2000); Najmaie et al. (2003); Stevens et al. (2003b), we have used a framework based on (i) a separation of the initial carrier photoinjection and the subsequent carrier scattering, and (ii) a perturbative expansion in the optical field amplitudes, with injection rates obtained in a Fermi’s golden rule limit for the bichromatic field. The injection rates for carrier population control, spin control, current injection, and spin current injection, have been described phenomenologically by tensors $`\xi _{(I)}`$, $`\zeta _{(I)}`$, $`\eta _{(I)}`$, and $`\mu _{(I)}`$, respectively Atanasov et al. (1996); Fraser et al. (1999); Najmaie et al. (2003); Stevens et al. (2003b); Stevens et al. (2004). Like previous theories, we have used the long-wavelength limit, and neglect nonlocal corrections to the interaction Hamiltonian. But whereas previous theories of ‘1+2’ photoinjection used the independent particle approximation, we have included excitonic effects. We have shown that excitonic effects cause (i) an enhancement of each ‘1+2’ process, and (ii) a phase shift for current injection and spin current injection. Our main results, the modifications of the aforementioned tensors relative to the independent particle approximation are given in (46), (51), (53), and (54). These particularly simple results are valid at photon energies for which transitions from the heavy- and light-hole bands are dominant; more general results are given for $`\eta _{(I)}`$ and $`\xi _{(I)}`$ in (38) and (48). Our results are based on the effective mass model of Wannier excitons; degenerate bands are included, but we use a spherical approximation to the exciton Hamiltonian, and we neglect envelope-hole coupling. This is a good approximation for many typical semiconductors, including GaAs, since the electron-hole envelope function extends over many unit cells due to the screening of the Coulomb interaction by the static dielectric constant Baldereschi and Lipari (1970, 1971, 1973); Sondergeld (1977b, a). As a consequence of making the spherical approximation, the phase shifts and Coulomb enhancements we find in this paper are independent of crystal orientation. Also, our results are limited to low excess photon energy since (i) the Wannier exciton Hamiltonian assumes parabolic Bloch bands, and (ii) we have truncated the expansion in $`𝐤`$ of the Bloch state velocity matrix elements, which is the basis of the transition amplitude expansion. By comparing the black dashed line and grey dash-dotted line in Fig. 1, one sees that higher order terms in $`𝐤`$ (for both bands and velocities) are important in GaAs for excess photon energies greater than about 200 meV. This can then be considered the limit of validity of our calculation. However, combining the Coulomb enhancement calculated assuming parabolic bands with the nonperturbative independent particle approximation result (as was done for the solid black line in Fig. 1) likely gives a good approximation for a few hundred more meV; this was the case for one- and two-photon absorption Sturge (1962); Weiler (1981). It is interesting to ask if there are other sources of intrinsic phases to the current (or spin current) besides the one that we have identified here, as these may produce spectral features in the intrinsic phase. One possibility is the coupling between bound $`so`$-$`c`$ excitons and the unbound $`hh`$-$`c`$ or $`lh`$-$`c`$ excitons, since it is known that the intrinsic phase can show spectral features near a resonance Seideman (1999). Another possibility is the envelope-hole coupling between the continua of unbound $`hh`$-$`c`$ and $`lh`$-$`c`$ excitons that was neglected in our treatment. The effect of coupled continua in the general theory of ‘$`n+m`$’ phase shifts has not been considered to date. Finally, we note that the intrinsic phase and Coulomb enhancement may be greater in reduced dimensional systems, which have greater exciton binding energies. The carrier-carrier Coulomb interaction was included in the theory for ‘$`1+2`$’ control of electrons in biased asymmetric quantum wells, although the intrinsic phase was not studied Pötz (1998). ## Appendix A Intermediate state Coulomb enhancement Consider the functions $`N^{(\text{a-f})}`$ and $`N^{(\text{a-a})}`$, which appear in (34) and (36); we refer to them collectively as $`N`$. First, note that due to the energy conserving delta function in (7), $`\kappa `$ will be equal to $`\kappa _{cv}`$ \[see Eq. (39)\], and thus $`N`$ is a function only of $`\omega `$, $`E_{cv}^g`$, $`E_{c^{}v^{}}^g`$, $`\mu _{cv}`$, and $`\mu _{c^{}v^{}}`$. Second, note that $`N`$ is defined so that if the electron-hole attraction is turned off, for example by letting $`ϵ\mathrm{}`$, then $`N1`$ foo (c). This allows $`N`$ to be identified as part of the Coulomb enhancement. In particular, $`N`$ is the enhancement due to the Coulomb interaction in the intermediate states; if the Coulomb interaction is neglected for the intermediate states, $`N=1`$ Doni et al. (1974). \[Note that Lee and Fan Lee and Fan (1974) did not allow for $`v^{}v`$ in $`N`$ (related to $`J_j`$ in their notation).\] Since the integrand is smooth for the parameter range of interest, numerical integration of $`N`$ is straightforward; however, it need not be undertaken. Further simplification is possible since the parameter $`\gamma `$ can be considered to be much less than one. Since most materials have an exciton binding energy that is much smaller than the band gap, $`\mathrm{}\omega `$ is detuned from the band edge by many exciton binding energies at photon energies consistent with the approximations made here. In GaAs, for example, when $`2\mathrm{}\omega `$ is within 500 meV of the gap, $`\gamma `$ is at most $`0.09`$. An expansion of $`N^{(\text{a-f})}`$ for small $`\gamma `$, $$N^{(\text{a-f})}=1+\frac{2}{3}\gamma _{c^{}v^{}}+\left(\frac{4}{3}\mathrm{ln}2\frac{1}{3}\right)\gamma _{c^{}v^{}}^2+\left(S_0\frac{2}{15}a_{cv}^2\kappa ^2\right)\gamma _{c^{}v^{}}^3+O\left(\gamma _{c^{}v^{}}^4\right),$$ where $`S_00.5633`$, shows that $`N^{(\text{a-f})}`$ is approximately $`1`$ and nearly constant as a function of $`\omega `$. The same is true of $`N^{(\text{a-a})}`$, which has the expansion $$N^{(\text{a-a})}=1\frac{2}{a_{cv}}\left(a_{c^{}v^{}}a_{cv}\right)P+O\left(\gamma _{c^{}v^{}}^4\right),$$ where, with $`S_11.645`$, $$P\gamma _{c^{}v^{}}+\left(2\mathrm{ln}2\frac{a_{c^{}v^{}}}{a_{cv}}\right)\gamma _{c^{}v^{}}^2+\left(\frac{2a_{c^{}v^{}}^2}{3a_{cv}^2}\frac{2a_{c^{}v^{}}}{a_{cv}}\frac{1}{3}a_{c^{}v^{}}^2\kappa ^2+S_1\right)\gamma _{c^{}v^{}}^3.$$ In fact, when $`\mu _{cv}=\mu _{c^{}v^{}}`$, $`N^{(\text{a-a})}=1`$ even to fourth order in $`\gamma _{c^{}v^{}}`$. Fig. 4 shows a numerical integration of $`N^{(\text{a-f})}`$ using the parameters of GaAs. ## Appendix B Evaluation of the current injection tensor The tensor $`\eta _{cc^{}vv^{}}^{ijkl}`$, defined in (41), can be used to calculate the current injection tensor with or without excitonic effects using (38) or (44). It can be evaluated analytically in the approximation of parabolic bands. Part of the result for the 8 band Kane model has been given before but without the split-off band as an initial or intermediate state Bhat and Sipe (2000). We here give more detail, but only for the two-band terms. We denote the bands by a double index $`n`$ and $`s`$, where $`n`$ is one of $`\{c,hh,lh,so\}`$, and $`s`$ runs over the two spin states for each band. Since the Coulomb corrections to $`\eta _{(I)}`$ in (38) do not depend on the spin index, we can include the sum over spin indices from (38) in $`\eta _{c,v,v^{}}^{ijkl}`$. And, since we are only calculating two-band terms, we set $`c^{}=c`$ and $`v^{}=v`$. Thus $`\eta _{ccvv}^{ijkl}`$ $`=`$ $`i{\displaystyle \frac{\pi e^4}{\mathrm{}^2\omega ^3}}{\displaystyle \underset{s,s^{}}{}}{\displaystyle \frac{1}{L^3}}{\displaystyle \underset{𝐤}{}}\mathrm{\Delta }_{cv}^i{\displaystyle \frac{\{\left(𝐯_{cc}𝐯_{vv}\right),𝐯_{cs,vs^{}}^{}\}^{jk}}{E_{cv}\mathrm{}\omega }}v_{cs,vs^{}}^l\delta \left(2\omega E_{cv}\left(k\right)/\mathrm{}\right)`$ $`=`$ $`i{\displaystyle \frac{\pi e^4}{\mathrm{}^3\omega ^4}}{\displaystyle \frac{1}{8\pi ^3}}{\displaystyle }k^2dk{\displaystyle \frac{\mu _{cv}}{\mathrm{}k_{cv}}}\delta (kk_{cv})d\mathrm{\Omega }{\displaystyle \underset{s,s^{}}{}}{\displaystyle \frac{\mathrm{}k^i}{\mu _{cv}}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{}k^j}{\mu _{cv}}}v_{vs^{},cs}^kv_{cs,vs^{}}^l+(jk)`$ $`=`$ $`i{\displaystyle \frac{e^4}{\mathrm{}^2\omega ^4}}{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle \frac{k_{cv}^3}{\mu _{cv}}}{\displaystyle \frac{1}{2}}{\displaystyle }d\mathrm{\Omega }\widehat{k}^i\widehat{k}^j{\displaystyle \underset{s,s^{}}{}}v_{vs^{},cs}^kv_{cs,vs^{}}^l+(jk)`$ where $$k_{cv}=\sqrt{\frac{2\mu _{cv}}{\mathrm{}^2}\left(2\mathrm{}\omega E_{cv}^g\right)}.$$ For $`v=lh`$ or $`hh`$, $`E_{cv}^g=E_g`$, while for $`v=so`$, $`E_{cv}^g=E_g+\mathrm{\Delta }`$, where $`E_g`$ is the fundamental band gap, and $`\mathrm{\Delta }`$ is the spin-orbit splitting. The interband velocity matrix elements are approximated by their value at $`k=0`$, but still depend on the direction of $`𝐤`$. In terms of the orthogonal triple of unit vectors $`\widehat{𝐤}`$, $`\widehat{𝐥}`$ and $`\widehat{𝐦}`$, $`\widehat{𝐤}`$ $`=`$ $`\mathrm{sin}\theta \mathrm{cos}\varphi \widehat{𝐱}+\mathrm{sin}\theta \mathrm{sin}\varphi \widehat{𝐲}+\mathrm{cos}\theta \widehat{𝐳}`$ $`\widehat{𝐥}`$ $`=`$ $`\mathrm{cos}\theta \mathrm{cos}\varphi \widehat{𝐱}+\mathrm{cos}\theta \mathrm{sin}\varphi \widehat{𝐲}\mathrm{sin}\theta \widehat{𝐳}`$ $`\widehat{𝐦}`$ $`=`$ $`\mathrm{sin}\varphi \widehat{𝐱}+\mathrm{cos}\varphi \widehat{𝐲},`$ these matrix elements are: $`𝐯_{c,s;hh,s^{}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{E_P}{m}}}\left[\widehat{𝐥}𝕀+i\widehat{𝐦}\sigma _z\right]_{s,s^{}}`$ $`𝐯_{c,s;lh,s^{}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{E_P}{3m}}}\left[2\widehat{𝐤}𝕀+i\widehat{𝐥}\sigma _yi\widehat{𝐦}\sigma _x\right]_{s,s^{}}`$ $`𝐯_{c,s;so,s^{}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{E_P}{6m}}}\left[\widehat{𝐤}𝕀i\widehat{𝐥}\sigma _y+i\widehat{𝐦}\sigma _x\right]_{s,s^{}},`$ where $`E_P`$ is the Kane energy Kane (1957). Here, $`𝕀`$ is the $`2\times 2`$ identity matrix and $`\sigma _i`$ are the Pauli spin matrices. Of course, for parabolic bands, the intraband matrix elements are $`n,s,𝐤\left\|𝐯\right\|n,s^{},𝐤=\delta _{s,s^{}}\widehat{𝐤}\mathrm{}k/m_n`$, where $`m_n`$ is the effective mass of band $`n`$. (In the proper Kane model, the effective masses are given in terms of the parameters $`E_g`$, $`\mathrm{\Delta }`$, and $`E_P`$, but we treat them as additional parameters, which is equivalent to including remote band effects on the effective masses.) The sums over spin then yield $`{\displaystyle \underset{s,s^{}}{}}v_{hh,s^{};c,s}^kv_{c,s;hh,s^{}}^l`$ $`=`$ $`{\displaystyle \frac{E_P}{2m}}\left(\delta _{k,l}\widehat{k}^k\widehat{k}^l\right)`$ $`{\displaystyle \underset{s,s^{}}{}}v_{lh,s^{};c,s}^kv_{c,s;lh,s^{}}^l`$ $`=`$ $`{\displaystyle \frac{E_P}{2m}}\left(\widehat{k}^k\widehat{k}^l+{\displaystyle \frac{1}{3}}\delta _{k,l}\right)`$ $`{\displaystyle \underset{s,s^{}}{}}v_{so,s^{};c,s}^kv_{c,s;so,s^{}}^l`$ $`=`$ $`{\displaystyle \frac{E_P}{3m}}\delta _{k,l}.`$ The remaining angular integrals can be done using $`{\displaystyle 𝑑\mathrm{\Omega }\widehat{k}^i\widehat{k}^j}`$ $`=`$ $`{\displaystyle \frac{4\pi }{3}}\delta _{i,j}`$ $`{\displaystyle 𝑑\mathrm{\Omega }\widehat{k}^i\widehat{k}^j\widehat{k}^k\widehat{k}^l}`$ $`=`$ $`{\displaystyle \frac{4\pi }{15}}\left(\delta _{i,j}\delta _{k,l}+\delta _{i,k}\delta _{j,l}+\delta _{i,l}\delta _{j,k}\right).`$ The result for $`\eta _{(I\text{-free})}`$ is $$\eta _{(I\text{-free})}^{ijkl}=i\frac{\sqrt{2}}{9\pi }\frac{e^4E_P}{\omega ^4\mathrm{}^5\sqrt{m}}\left[\left(2\mathrm{}\omega E_g\right)^{\frac{3}{2}}\left(T_{hh}^{ijkl}+T_{lh}^{ijkl}\right)+\left(2\mathrm{}\omega E_g\mathrm{\Delta }\right)^{\frac{3}{2}}T_{so}^{ijkl}\right],$$ (63) where the tensor properties are $`T_{hh}^{ijkl}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mu _{c,hh}}{m}}}\left({\displaystyle \frac{9}{20}}\delta _{i,j}\delta _{k,l}+{\displaystyle \frac{9}{20}}\delta _{i,k}\delta _{j,l}{\displaystyle \frac{3}{10}}\delta _{i,l}\delta _{j,k}\right)`$ $`T_{lh}^{ijkl}`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mu _{c,lh}}{m}}}\left({\displaystyle \frac{11}{20}}\delta _{i,j}\delta _{k,l}+{\displaystyle \frac{11}{20}}\delta _{i,k}\delta _{j,l}+{\displaystyle \frac{3}{10}}\delta _{i,l}\delta _{j,k}\right)`$ $`T_{so}^{ijkl}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mu _{c,so}}{m}}}\left(\delta _{i,j}\delta _{k,l}+\delta _{i,k}\delta _{j,l}\right),`$ and the term involving $`T_{so}^{ijkl}`$ should not be included if $`2\mathrm{}\omega <E_g+\mathrm{\Delta }`$. The free particle current injection tensor has also been investigated for parabolic bands, but with a simple three-band model Sheik-Bahae (1999). That model does not have the matrix elements $`𝐯_{c,lh}`$ and $`𝐯_{c,lh}`$, and thus differs from our result for $`T_{lh}^{xxxx}`$. ###### Acknowledgements. This work was financially supported by the Natural Science and Engineering Research Council, Photonics Research Ontario, and the US Defense Advanced Research Projects Agency. We gratefully acknowledge many stimulating discussions with Daniel Côté, James Fraser, Ali Najmaie, Fred Nastos, Eugene Sherman, Art Smirl, Marty Stevens, and Henry van Driel.
warning/0506/math-ph0506055.html	ar5iv	text	# VARIATIONAL PROBLEMS IN ELASTIC THEORY OF BIOMEMBRANES, SMECTIC-A LIQUID CRYSTALS, AND CARBON RELATED STRUCTURES ## 1 Introduction The morphology of thin structures (always represented by a smooth surface $`M`$ in this paper) is an old problem. First, we look back on the history . As early as in 1803, Plateau studied a soap film attaching to a metallic ring when the ring passed through soap water . By taking the minimum of the free energy $`F=\lambda _MdA`$, he obtained $`H=0`$, where $`\lambda `$ and $`H`$ are the surface tension and mean curvature of the soap film, respectively. From 1805 and 1806, Young and Laplace studied soap bubbles. By taking the minimum of the free energy $`F=pdV+\lambda _MdA`$, they obtained $`H=p/2\lambda `$, where $`p`$ is the osmotic pressure (pressure difference between outer and inner sides) of a soap bubble and $`V`$ is the volume enclosed by the bubble. We can only observe spherical bubbles because “An embedded surface with constant mean curvature in 3-dimensional (3D) Euclidian space ($`𝔼^3`$) must be a spherical surface” . In 1812, Poisson considered a solid shell and put forward the free energy $`F=_MH^2dA`$. Its Euler-Lagrange equation is $`^2H+2H(H^2K)=0`$ . Now the solutions to this equation are called Willmore surfaces. In 1973, Helfrich recognized that lipid bilayers could be regarded as smectic-A (SmA) liquid crystals (LCs) at room temperature. Based on the elastic theory of liquid crystals , he proposed the curvature energy per unit area of the bilayer $`_{lb}=(k_c/2)(2H+c_0)^2+\overline{k}K,`$ (1) where $`k_c`$ and $`\overline{k}`$ are elastic constants. $`K`$ and $`c_0`$ are Gaussian curvature and spontaneous curvature of the lipid bilayer, respectively. Starting with Helfrich’s curvature energy (1), the morphology of lipid vesicles has been deeply understood . Especially, the free energy is expressed as $`F=pdV+_M(\lambda +_{lb})dA`$ for lipid vesicles, and the corresponding Euler-Lagrange equation is : $`p2\lambda H+k_c^2(2H)+k_c(2H+c_0)(2H^2c_0H2K)=0.`$ (2) For an open lipid bilayer with a free edge $`C`$, the free energy is expressed as $`F=_M(\lambda +_{lb})dA+\gamma _Cds`$, where $`\gamma `$ is the line tension of the edge. The corresponding Euler-Lagrange equations are as follows : $`k_c(2H+c_0)(2H^2c_0H2K)2\lambda H+k_c^2(2H)`$ $`=0`$ (3) $`\left[k_c(2H+c_0)+\overline{k}k_n\right]\|_C`$ $`=0`$ (4) $`\left[2k_c{\displaystyle \frac{H}{𝐞_2}}+\gamma k_n+\overline{k}{\displaystyle \frac{d\tau _g}{ds}}\right]\|_C`$ $`=0`$ (5) $`\left[{\displaystyle \frac{k_c}{2}}(2H+c_0)^2+\overline{k}K+\lambda +\gamma k_g\right]\|_C`$ $`=0,`$ (6) where $`k_n`$, $`k_g`$, and $`\tau _g`$ are normal curvature, geodesic curvature, and geodesic torsion of the boundary curve. The unit vector $`𝐞_2`$ (see also Fig. 1) is perpendicular to tangent vector of edge $`C`$ and normal vector of surface $`M`$. Above four equations are called the shape equation and boundary conditions of open lipid bilayers. The boundary conditions are available for open lipid bilayers with more than one edge because the edge in our derivation is a general one. Secondly, we turn to the puzzle about the formation of focal conic structures in SmA LCs. As we imagine, the configuration of minimum energy in SmA LCs is a flat layer structure. But Dupin cyclides are usually formed when LCs cool from isotropic phase to SmA phase in the experiment . Why the cyclides are preferred to other geometrical structures under the preservation of the interlayer spacing ? This phenomenon can be understood by the concept that the Gibbs free energy difference between isotropic and SmA phases must be balanced by the curvature elastic energy of SmA layers . The total free energy includes curvature energy, volume energy and surface energy. It is expressed formally as $`F=(H,K,t)dA`$, where $`t`$ is the thickness of the focal conic domain; $`H`$ and $`K`$ are mean curvature and Gaussian curvature of the inmost layer surface, respectively. The Euler-Lagrange equations corresponding to the free energy are as follows : $`{\displaystyle (/t)dA}`$ $`=0`$ (7) $`(^2/2+2H^2K)/H+(\stackrel{~}{}+2KH)/K2H`$ $`=0.`$ (8) Solving both equations can give good explanation to focal conic domains . The new operator $`\stackrel{~}{}`$ can be found in in the appendix of Ref. . Thirdly, let us see carbon related structures. There are three typical structures composed of carbon atoms: Buckyball (C<sub>60</sub>), single-walled carbon nanotube (SWNT), and carbon torus. In the continuum limit, we derive the curvature energy of single graphitic layer $`E=\left[\frac{1}{2}k_c(2H)^2+\overline{k}K\right]dA`$ from the lattice model , where $`k_c`$ and $`\overline{k}`$ are elastic constants. The total free energy of a graphite layer is $`F=\left[\frac{1}{2}k_c(2H)^2+\overline{k}K\right]dA+\lambda dA`$, where $`\lambda `$ is the surface energy per unit area for graphite. Please note that the surface energy per unit area for solid structures is not as a constant quantity as the surface tension for fluid membranes. The Euler-Lagrange equation corresponding to the free energy is $`^2H+2H(H^2K)\lambda H/k_c=0`$. C<sub>60</sub> and carbon torus can be understood with $`\lambda =0`$, while SWNT satisfies $`R^2=k_c/2\lambda `$, where $`R`$ is its radius. The rest of this paper is organized as follows. In Sec. 2, we show how to derive the Euler-Lagrange equation from the free energy functional by using exterior differential forms. The method is developed in Refs. and , which might be equivalent to the work by Griffiths in essence. But it is more convenient to apply in variational problems on 2D surface in $`𝔼^3`$. In Sec. 3, we give several analytic solutions to the shape equation of lipid vesicles, and to the shape equation and boundary conditions of open lipid bilayers. In Sec. 4, we discuss the elasticity and stability of cell membranes. A brief summary is given in the last section. ## 2 Variational problems on 2D surface Many variational problems are shown in Introduction. Here we deal with them by using exterior differential forms. Let us consider a surface $`M`$ with an edge $`C`$ as shown in Fig. 1. At every point P in the surface, we can choose a right-handed, orthonormal frame {$`𝐞_1,𝐞_2,𝐞_3`$} with $`𝐞_3`$ being the normal vector. For a point in curve $`C`$, $`𝐞_1`$ is the tangent vector of $`C`$ such that $`𝐞_2`$ points to the side that the surface is located in. The differential of the frame is expressed as $`\mathrm{d}𝐫`$ $`=\omega _1𝐞_1+\omega _2𝐞_2`$ (9) $`\mathrm{d}𝐞_i`$ $`=\omega _{ij}𝐞_j,(i=1,2,3)`$ (10) where $`\omega _1`$, $`\omega _2`$, $`\omega _{ij}=\omega _{ji}`$ ($`i,j=1,2,3`$) are 1-forms. Please note that the repeated subindex represents the summation from 1 to 3, unless otherwise specified in this paper. The structure equations of the surface are as follows: $`\mathrm{d}\omega _1`$ $`=\omega _{12}\omega _2`$ (11) $`\mathrm{d}\omega _2`$ $`=\omega _{21}\omega _1`$ (12) $`\omega _{13}`$ $`=a\omega _1+b\omega _2,\omega _{23}=b\omega _1+c\omega _2`$ (13) $`\mathrm{d}\omega _{ij}`$ $`=\omega _{ik}\omega _{kj}(i,j=1,2,3),`$ (14) where $`a`$, $`b`$, $`c`$ are related to the curvatures with $`2H=a+c`$ and $`K=acb^2`$. The variation of the frame is denoted by $`\delta 𝐫`$ $`=\delta _1𝐫+\delta _2𝐫+\delta _3𝐫`$ (15) $`\delta _i𝐫`$ $`=\mathrm{\Omega }_i𝐞_i(i=1,2,3)`$ (16) $`\delta _l𝐞_i`$ $`=\mathrm{\Omega }_{lij}𝐞_j,(i,l=1,2,3),`$ (17) with $`\mathrm{\Omega }_{lij}=\mathrm{\Omega }_{lji}`$ ($`i,j,l=1,2,3`$). In equation (16), the repeated subindex does not represent summation. It is easy to prove that the operator $`\delta _l(l=1,2,3)`$ has the similar properties with the partial differential while the operator $`\delta `$ has the similar properties with the total differential operator . Using $`\mathrm{d}\delta _l𝐫=\delta _l\mathrm{d}𝐫`$ and $`\mathrm{d}\delta _l𝐞_i=\delta _l\mathrm{d}𝐞_i`$, we obtain variational equations of the frame as follows : $`\delta _1\omega _1`$ $`=\mathrm{d}\mathrm{\Omega }_1\omega _2\mathrm{\Omega }_{121}`$ (18) $`\delta _1\omega _2`$ $`=\mathrm{\Omega }_1\omega _{12}\omega _1\mathrm{\Omega }_{112}`$ (19) $`\mathrm{\Omega }_{113}`$ $`=a\mathrm{\Omega }_1,\mathrm{\Omega }_{123}=b\mathrm{\Omega }_1;`$ (20) $`\delta _2\omega _1`$ $`=\mathrm{\Omega }_2\omega _{21}\omega _2\mathrm{\Omega }_{221}`$ (21) $`\delta _2\omega _2`$ $`=\mathrm{d}\mathrm{\Omega }_2\omega _1\mathrm{\Omega }_{212}`$ (22) $`\mathrm{\Omega }_{213}`$ $`=b\mathrm{\Omega }_2,\mathrm{\Omega }_{223}=c\mathrm{\Omega }_2;`$ (23) $`\delta _3\omega _1=\mathrm{\Omega }_3\omega _{31}\omega _2\mathrm{\Omega }_{321}`$ (24) $`\delta _3\omega _2=\mathrm{\Omega }_3\omega _{32}\omega _1\mathrm{\Omega }_{312}`$ (25) $`\mathrm{d}\mathrm{\Omega }_3=\mathrm{\Omega }_{313}\omega _1+\mathrm{\Omega }_{323}\omega _2;`$ (26) $`\delta _l\omega _{ij}=\mathrm{d}\mathrm{\Omega }_{lij}+\mathrm{\Omega }_{lik}\omega _{kj}\omega _{ik}\mathrm{\Omega }_{lkj}.`$ (27) Using them, we can prove that $`\delta _1(\omega _1\omega _2)`$ $`=\mathrm{d}(\omega _2\mathrm{\Omega }_1)`$ (28) $`\delta _2(\omega _1\omega _2)`$ $`=\mathrm{d}(\omega _1\mathrm{\Omega }_2)`$ (29) $`\delta _3(\omega _1\omega _2)`$ $`=2H\mathrm{\Omega }_3\omega _1\omega _2`$ (30) $`\delta _3(2H)\omega _1\omega _2`$ $`=2(2H^2K)\mathrm{\Omega }_3\omega _1\omega _2+\mathrm{d}\mathrm{d}\mathrm{\Omega }_3`$ (31) $`\delta _3K\omega _1\omega _2`$ $`=2KH\mathrm{\Omega }_3\omega _1\omega _2+\mathrm{d}\stackrel{~}{}\stackrel{~}{\mathrm{d}}\mathrm{\Omega }_3,`$ (32) where $``$ is a function of $`2H`$ and $`K`$. $``$ is Hodge star operator satisfying the following properties: (i) $`f=f\omega _1\omega _2`$ for scalar function $`f`$; (ii) $`\omega _1=\omega _2`$, $`\omega _2=\omega _1`$. $`\stackrel{~}{}`$ and $`\stackrel{~}{\mathrm{d}}`$ are new operators defined in Ref. that satisfy: (i) If $`df=f_1\omega _1+f_2\omega _2`$, then $`\stackrel{~}{\mathrm{d}}f=f_1\omega _{13}+f_2\omega _{23}`$; (ii) $`\stackrel{~}{}\omega _{13}=\omega _{23},\stackrel{~}{}\omega _{23}=\omega _{13}`$; (iii) $`_M(f\mathrm{d}\stackrel{~}{\mathrm{d}}hh\mathrm{d}\stackrel{~}{\mathrm{d}}f)=_M(f\stackrel{~}{\mathrm{d}}hh\stackrel{~}{\mathrm{d}}f)`$, $`_M(f\mathrm{d}\stackrel{~}{}\stackrel{~}{\mathrm{d}}hh\mathrm{d}\stackrel{~}{}\stackrel{~}{\mathrm{d}}f)=_M(f\stackrel{~}{}\stackrel{~}{\mathrm{d}}hh\stackrel{~}{}\stackrel{~}{\mathrm{d}}f)`$ for any smooth functions $`f`$ and $`h`$ on $`M`$. <sup>1</sup><sup>1</sup>1These two expressions are very similar to the second Green identity. Their proof can be found in Lemma 2.1 of Ref. . Now, we consider the variational problem on closed surface. In this case, the general functional is expressed as: $`F={\displaystyle _M}(2H,K)dA+p{\displaystyle _V}dV.`$ (33) Using Stokes theorem and the variational equations of the frame, we can prove that $`\delta _1{\displaystyle _V}dV`$ $`=\delta _2{\displaystyle _V}dV=0`$ (34) $`\delta _3{\displaystyle _V}dV`$ $`={\displaystyle _M}\mathrm{\Omega }_3dA.`$ (35) Combining them with equations (28)–(32), we have $`\delta _1F=\delta _2F=0`$, and the Euler-Lagrange equation corresponding to functional (33): $`\left[\left(^2+4H^22K\right){\displaystyle \frac{}{(2H)}}+\left(\stackrel{~}{}+2KH\right){\displaystyle \frac{}{K}}2H\right]+p=0.`$ (36) Next, we consider the variational problem on open surface with an edge $`C`$. In this case, the general functional is expressed as: $`F={\displaystyle _M}(2H,K)dA+{\displaystyle _C}\mathrm{\Gamma }(k_n,k_g)ds.`$ (37) Similarly, we derive its Euler-Lagrange equations as $`(^2+4H^22K){\displaystyle \frac{}{(2H)}}+(\stackrel{~}{}+2KH){\displaystyle \frac{}{K}}2H=0`$ (38) $`𝐞_2\left[{\displaystyle \frac{}{(2H)}}\right]+𝐞_2\stackrel{~}{}\left({\displaystyle \frac{}{K}}\right){\displaystyle \frac{d}{ds}}\left(\tau _g{\displaystyle \frac{}{K}}\right)+{\displaystyle \frac{d^2}{ds^2}}\left({\displaystyle \frac{\mathrm{\Gamma }}{k_n}}\right)+{\displaystyle \frac{\mathrm{\Gamma }}{k_n}}(k_n^2\tau _g^2)`$ $`+\tau _g{\displaystyle \frac{d}{ds}}\left({\displaystyle \frac{\mathrm{\Gamma }}{k_g}}\right)+{\displaystyle \frac{d}{ds}}\left(\tau _g{\displaystyle \frac{\mathrm{\Gamma }}{k_g}}\right)\left(\mathrm{\Gamma }{\displaystyle \frac{\mathrm{\Gamma }}{k_g}}k_g\right)k_n\|_C=0`$ (39) $`{\displaystyle \frac{}{(2H)}}k_n{\displaystyle \frac{}{K}}+{\displaystyle \frac{\mathrm{\Gamma }}{k_g}}k_n{\displaystyle \frac{\mathrm{\Gamma }}{k_n}}k_g\|_C=0`$ (40) $`{\displaystyle \frac{d^2}{ds^2}}\left({\displaystyle \frac{\mathrm{\Gamma }}{k_g}}\right)+K{\displaystyle \frac{\mathrm{\Gamma }}{k_g}}k_g\left(\mathrm{\Gamma }{\displaystyle \frac{\mathrm{\Gamma }}{k_g}}k_g\right)+2(k_nH)k_g{\displaystyle \frac{\mathrm{\Gamma }}{k_n}}`$ $`\tau _g{\displaystyle \frac{d}{ds}}\left({\displaystyle \frac{\mathrm{\Gamma }}{k_n}}\right){\displaystyle \frac{d}{ds}}\left(\tau _g{\displaystyle \frac{\mathrm{\Gamma }}{k_n}}\right)\|_C=0.`$ (41) In special cases, above equations (36), (38)–(41) are degenerated into different equations mentioned in Introduction. For example, if taking $`=_{lb}+\lambda `$ and $`\mathrm{\Gamma }=\gamma `$, These equations are degenerated into equations (3)–(6), respectively. ## 3 Morphology of lipid bilayers There are three typical solutions to shape equation (2): Sphere, torus, and biconcave vesicle. First, a spherical vesicle with radius $`R`$ satisfies equation (2) if $`pR^2+2\lambda Rk_cc_0(2c_0R)=0`$ is valid. Secondly, Ou-Yang used equation (2) to predict that a torus with the radii of two generating circles $`r`$ and $`R`$ satisfying $`R/r=\sqrt{2}`$ should be observed in lipid systems . This striking prediction has been confirmed experimentally by three groups . Thirdly, the first exact axisymmetric solution with biconcave shape, as shown in Fig. 2, was found under the condition of $`p=\lambda =0`$ as : $`z`$ $`=z_0+{\displaystyle _0^\rho }\mathrm{tan}\psi (\rho ^{})𝑑\rho ^{}`$ (42) $`\mathrm{sin}\psi (\rho )`$ $`=c_0\rho \mathrm{ln}(\rho /\rho _B),c_0>0,`$ (43) where $`z(\rho )`$ is the contour of the cross-section. $`z`$ axis is the rotational axis, and $`\psi (\rho )`$ the tangent angle of the contour at distance $`\rho `$. This solution can explain the classic physiological puzzle : Why the red blood cells of humans are always in biconcave shape? To the shape equation (3) and boundary conditions (4)–(6) of open lipid bilayers, we can find two analytical solutions : One is the central part of a torus and another is a cup-like membrane shown in Fig. 3. Numerical method and solutions to these equations can be find in Ref. . ## 4 Elasticity and stability of cell membranes A cell membrane is simplified as lipid bilayer plus membrane skeleton. The skeleton is a cross-linking protein network and joint to the bilayer at some points. We know that the cross-linking polymer structure also exists in rubber at molecular levels. Thus we can transplant the theory of rubber elasticity to describe the membrane skeleton. Based on Helfrich’s theory and physics of rubber elasticity, the free energy of a closed cell membrane can be expressed as : $`F={\displaystyle _M}(_d+_H)dA+p{\displaystyle _V}dV,`$ (44) with $`_H=(k_c/2)(2H+c_0)^2+\lambda `$ and $`_d=(k_d/2)[(2J)^2Q]`$, where $`2J=\epsilon _{11}+\epsilon _{22}`$, $`Q=\epsilon _{11}\epsilon _{22}\epsilon _{12}^2`$. Here $`\epsilon _{ij}`$ ($`i,j=1,2`$) represents the in-plane strain of the membrane. $`k_d`$ is the elastic constant representing the entropic elasticity of membrane skeleton. Using the method in Sec. 2, we obtain the shape equation and the in-plane strain equations of the cell membrane as : $`p2H(\lambda +k_dJ)+k_c(2H+c_0)(2H^2c_0H2K)`$ $`+k_c^2(2H){\displaystyle \frac{k_d}{2}}(a\epsilon _{11}+2b\epsilon _{12}+c\epsilon _{22})=0`$ (45) $`k_d[\mathrm{d}(2J)\omega _2{\displaystyle \frac{1}{2}}(\epsilon _{11}\mathrm{d}\omega _2\epsilon _{12}\mathrm{d}\omega _1)+{\displaystyle \frac{1}{2}}\mathrm{d}(\epsilon _{12}\omega _1+\epsilon _{22}\omega _2)]=0`$ (46) $`k_d[\mathrm{d}(2J)\omega _1{\displaystyle \frac{1}{2}}(\epsilon _{12}\mathrm{d}\omega _2\epsilon _{22}\mathrm{d}\omega _1){\displaystyle \frac{1}{2}}\mathrm{d}(\epsilon _{11}\omega _1+\epsilon _{12}\omega _2)]=0.`$ (47) An obvious solution is the spherical cell membrane with homogenous strains: $`\epsilon _{11}=\epsilon _{22}=\epsilon `$ (a constant) and $`\epsilon _{12}=0`$. The radius $`R`$ of the sphere must satisfy $`pR^2+(2\lambda +3k_d\epsilon )R+k_cc_0(c_0R2)=0.`$ (48) Now we will show the biological function of membrane skeleton by discussing the mechanical stability of a spherical cell membrane. Using Hodge decomposed theorem , $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ can be expressed as $`\mathrm{\Omega }_1\omega _1+\mathrm{\Omega }_2\omega _2=d\mathrm{\Omega }+d\chi `$ by two scalar functions $`\mathrm{\Omega }`$ and $`\chi `$. Through complex calculations, we obtain the second order variation of the free energy for spherical membrane $`\delta ^2=G_1+G_2`$, where $`G_1`$ $`={\displaystyle _M}\mathrm{\Omega }_3^2\{3k_d/R^2+(2k_cc_0/R^3)+p/R\}dA`$ $`+{\displaystyle _M}\mathrm{\Omega }_3^2\mathrm{\Omega }_3\{k_cc_0/R+2k_c/R^2+pR/2\}dA`$ $`+{\displaystyle _M}k_c(^2\mathrm{\Omega }_3)^2dA+{\displaystyle \frac{3k_d}{R}}{\displaystyle _M}\mathrm{\Omega }_3^2\mathrm{\Omega }\mathrm{d}A`$ $`+k_d{\displaystyle _M}\left(^2\mathrm{\Omega }\right)^2dA+{\displaystyle \frac{k_d}{2R^2}}{\displaystyle _M}\mathrm{\Omega }^2\mathrm{\Omega }\mathrm{d}A,`$ (49) and $`G_2=(k_d/4)_M^2\chi (^2\chi +2\chi /R^2)dA`$. Because $`G_2`$ is positive definite, we merely need to discuss $`G_1`$. $`\mathrm{\Omega }_3`$ and $`\mathrm{\Omega }`$ in the expression of $`G_1`$ are arbitrary functions defined in a sphere and can be expanded by spherical harmonic functions : $`\mathrm{\Omega }_3=_{l=0}^{\mathrm{}}_{m=l}^{m=l}a_{lm}Y_{lm}(\theta ,\varphi )`$ and $`\mathrm{\Omega }=_{l=0}^{\mathrm{}}_{m=l}^{m=l}b_{lm}Y_{lm}(\theta ,\varphi )`$ with $`a_{lm}^{}=(1)^ma_{l,m}`$ and $`b_{lm}^{}=(1)^mb_{l,m}`$. Considering (48), we write $`G_1`$ in a quadratic form: $`G_1`$ $`={\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{l}{}}}2\|a_{lm}\|^2\{3k_d+[l(l+1)2][l(l+1)k_c/R^2k_cc_0/RpR/2]\}`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{l}{}}}{\displaystyle \frac{3k_d}{R}}l(l+1)(a_{lm}^{}b_{lm}+a_{lm}b_{lm}^{})`$ $`+{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{l}{}}}{\displaystyle \frac{k_d}{R^2}}\left[2l^2(l+1)^2l(l+1)\right]\|b_{lm}\|^2.`$ (50) It is easy to prove that, if $`p<p_l=\frac{3k_d}{\left[2l(l+1)1\right]R}+\frac{2k_c[l(l+1)c_0R]}{R^3}(l=2,3,\mathrm{})`$, then $`G_1`$ is positive definite. Thus we must take the minimum of $`p_l`$ to obtain the critical pressure <sup>2</sup><sup>2</sup>2In Ref. , we ignore the effect of in-plane modes $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ on the critical pressure and obtain the invalid value.: $`p_c=\mathrm{min}\{p_l\}=\{\begin{array}{c}\frac{3k_d}{11R}+\frac{2k_c[6c_0R]}{R^3}<\frac{k_c[232c_0R]}{R^3},(3k_dR^2<121k_c)\\ \frac{2\sqrt{3k_dk_c}}{R^2}+\frac{k_c}{R^3}(12c_0R),(3k_dR^2>121k_c).\end{array}`$ (53) Taking typical data of cell membrane, $`k_c20k_BT`$ , $`k_d6\times 10^4k_BT/nm^2`$ , $`h4nm`$, $`R1\mu m`$, $`c_0R1`$, we have $`p_c2`$ Pa from equation (53). If not considering membrane skeleton, that is $`k_d=0`$, we obtain $`p_c0.2`$ Pa. Therefore, membrane skeleton enhances the mechanical stability of cell membranes, at least for spherical shape. ## 5 Summary In above discussion, we introduce several problems in the elasticity of biomembranes, smectic-A liquid crystal, and carbon related structures. We deal with these variational problems on 2D surface by using exterior differential forms. Elasticity and stability of lipid bilayers and cell membranes are calculated and compared with each other. It is shown that membrane skeleton enhances the mechanical stability of cell membranes. ## Acknowledgement ZCT would like to thank the useful discussion and kind help of Prof. Y. S. Cho, T. Ivey, I. Mladenov, V. M. Vassilev, O. Yampolsky, and I. Zlatanov during this conference.
warning/0506/cond-mat0506496.html	ar5iv	text	# Effects of trap anisotropy on impurity scattering regime in a Fermi gas ## I Introduction Experiments on ballistic expansion and collective-mode excitations have provided important diagnostic tools in the study of quantum-degenerate atomic and molecular gases (for a recent review see Minguzzi et al. (2004)). A milestone example is the interpretation of the anisotropic expansion of a cloud of ultracold bosons as the smoking gun of Bose-Einstein condensation in a trapped dilute gas Anderson et al. (1995); Bradley et al. (1995); Davis et al. (1995). The later achievement of degenerate Fermi gases DeMarco and Jin (1999) and boson-fermion mixtures Truscott et al. (2001); Schreck et al. (2001); Goldwin et al. (2002) has opened up the possibility to experimentally observe rich quantum phase diagrams and to study the interplay between external confinement and interspecies collisions. More recently, the expansion behavior of strongly interacting fermion clouds has been the object of experiments aimed at probing novel superfluid states on the crossover from a BCS superfluid to a condensate of molecular dimers O’Hara et al. (2002); Greiner et al. (2003); Chin et al. (2004). Collective modes provide a direct measure of the collisionality of a dilute quantum gas. Monopolar and quadrupolar modes have been proposed as markers of the approach of quantum phase transitions in binary mixtures of Bose-Einstein condensed gases (BEC’s) Graham and Walls (1998), in boson-fermion mixtures Capuzzi et al. (2003a, b), and in fermion mixtures across the BCS-BEC crossover Kinast et al. (2004a); Bartenstein et al. (2004); Kinast et al. (2004b). The dynamical transition from collisionless to hydrodynamic behavior in fermion mixtures has been followed experimentally Gensemer and Jin (2001) and numerically Toschi et al. (2003) in two-component fermion mixtures by studying their dipolar oscillation modes. The role played in this context by mobile impurities inside a fermion gas under spherical confinement has also been studied by numerical means Capuzzi et al. (2004). In the present work we examine how the anisotropy of the trap affects the low-lying oscillation modes and the ballistic expansion of a gas of fermionic <sup>40</sup>K atoms containing a small concentration of thermal <sup>87</sup>Rb atoms. We numerically solve the Vlasov-Landau equations for the evolution of the phase-space distribution functions within a particle-in-cell approach and compare the results with simple scaling Ansatzes. Our results demonstrate that in a cigar-shaped harmonic confinement, where two different confinement frequencies determine two different time scales, collisionless and hydrodynamic behaviors can coexist in the low-lying collective excitations of the gas and that the aspect ratio of the expanding cloud shows a non-monotonic dependence on the anisotropy of the trap. The paper is organized as follows. In Sect. II we introduce the mixture under study and the basic equations that will be used to describe its dynamics. Section III analyzes the monopolar, dipolar, and quadrupolar oscillations of the gas as functions of the anisotropy of the confining potential, while Sect. IV gives a discussion of the dynamics of the free expansion of the fermion cloud. Finally, Sect. V presents a summary and the main conclusions of our work. ## II The mixture We consider a gas of fermionic atoms of mass $`m_F`$ confined inside an axially symmetric harmonic trap of the form $$V_F(𝐫)=\frac{1}{2}m_F\omega _{F,}^2\left(x^2+y^2+\lambda ^2z^2\right),$$ (1) where $`\omega _{F,}`$ is the angular trap frequency for motions along the $`\widehat{x}`$ and $`\widehat{y}`$ directions and $`\lambda =\omega _{F,z}/\omega _{F,}`$ is the anisotropy parameter. Potentials with $`\lambda 1`$ generate cigar-shaped density profiles. The fermions are wholly spin-polarized and hence at very low temperature the Pauli principle quenches collisions among them. A second component must then be added in order to thermalize the gas and to drive its collisionality. We dope the fermion gas with a small number of bosons having larger atomic mass $`m_B`$. We simply think of them as impurities confined by the harmonic potential $$V_B(𝐫)=\frac{1}{2}m_B\omega _{B,}^2\left(x^2+y^2+\lambda ^2z^2\right),$$ (2) where $`\omega _{B,}`$ is the angular trap frequency in the azimuthal plane and we have assumed that bosons and fermions share the same trap anisotropy $`\lambda `$. The difference in atomic masses and particle numbers causes a much lower quantum degeneracy temperature for the bosons and we can therefore consider them as an uncondensed cloud, even in the presence of highly degenerate fermions. The dynamics of the mixture is described through the one-body distribution functions $`f^{(F,B)}(𝐫,𝐩,t)`$ in the Boltzmann approximation. Their evolution is governed by the Vlasov-Landau kinetic equations (VLE), $$_tf^{(j)}+\frac{𝐩}{m_j}_𝐫f^{(j)}_𝐫U^{(j)}_𝐩f^{(j)}=C[f^{(F)},f^{(B)}]$$ (3) where the Hartree-Fock effective potential is $`U^{(j)}(𝐫,t)=V_j(𝐫)+gn^{(\overline{j})}(𝐫,t)`$ with $`\overline{j}`$ denoting the species different from $`j`$. Here we have set $`g=2\pi \mathrm{}^2a/m_r`$ with $`a`$ being the $`s`$-wave scattering length of a fermion-boson pair and $`m_r`$ its reduced mass, and $`n^{(j)}(𝐫,t)`$ is the spatial density given by integration of $`f^{(j)}(𝐫,𝐩,t)`$ over the momentum degrees of freedom. Since we deal with low concentrations of impurities we have neglected impurity-impurity interactions. In addition, collisions between spin-polarized fermions are negligible at low temperature and thus the collision integral $`C`$ in Eq. (3) involves only collisions between fermions and impurities. This is given by $`C`$ $`=`$ $`{\displaystyle \frac{\sigma }{4\pi (2\pi \mathrm{})^3}}{\displaystyle }d^3p_2d\mathrm{\Omega }_fv[(1f^{(F)})(1+f_2^{(B)})f_3^{(F)}f_4^{(B)}`$ (4) $`f^{(F)}f_2^{(B)}(1f_3^{(F)})(1+f_4^{(B)})],`$ where $`f^{(j)}f^{(j)}(𝐫,𝐩,t)`$ and $`f_i^{(j)}f^{(j)}(𝐫,𝐩_i,t)`$, $`d\mathrm{\Omega }_f`$ is the element of solid angle for the outgoing relative momentum $`𝐩_3𝐩_4`$, $`v=\|𝐯𝐯_2\|`$ is the relative velocity of the incoming particles, and $`\sigma =4\pi a^2`$ is the scattering cross-section. The collision satisfies conservation of momentum ($`𝐩+𝐩_2=𝐩_3+𝐩_4`$) and energy ($`\epsilon +\epsilon _2=\epsilon _3+\epsilon _4`$), with $`\epsilon _j=p_j^2/2m_j+U^{(j)}`$. The solution of Eqs. (3) is carried out by using a numerical approach based on particle-in-cell plus Monte Carlo sampling techniques, which allows us to evaluate the kinetics of such systems down to $`T0.1T_F`$. The technical details of the method have been given elsewhere Toschi et al. (2003); Capuzzi et al. (2004); Succi et al. (2004). In the following we shall focus on the dependence of the dynamics of a mixture of $`N_F=10^4`$ <sup>40</sup>K atoms and $`N_B=10^2`$ <sup>87</sup>Rb atoms on the anisotropy $`\lambda `$ of the traps at fixed values of the average trap frequencies $`\overline{\omega }_F=(\omega _{F,}^2\omega _{F,z})^{1/3}`$ and $`\overline{\omega }_B=(\omega _{B,}^2\omega _{B,z})^{1/3}`$. These are taken as the geometric averages of the trap frequencies in the experiments carried out at LENS on <sup>40</sup>K-<sup>87</sup>Rb mixtures Ferlaino et al. (2003), i.e. we set $`\overline{\omega }_F=2\pi \times 134.1`$ s<sup>-1</sup> and $`\overline{\omega }_B=2\pi \times 91.2`$ s<sup>-1</sup>. We fix the temperature at $`T=0.2T_F`$, with $`T_F=\mathrm{}\overline{\omega }_F(6N_F)^{1/3}/k_B`$ being the Fermi degeneracy temperature for noninteracting fermions. This also corresponds to $`T2.7T_{\text{BEC}}`$ where $`T_{\text{BEC}}=0.94\mathrm{}\overline{\omega }_B(N_B)^{1/3}/k_B`$ is the condensation temperature for the noninteracting Bose component. Finally, we assume a repulsive $`s`$-wave scattering length $`a=2000`$ Bohr radii. This choice of the fermion-boson scattering length is dictated by computational convenience, but still leaves the mixture at $`\lambda =1`$ in a collisionless-to-intermediate scattering regime as for the <sup>40</sup>K-<sup>87</sup>Rb mixtures studied at LENS. ## III Oscillation modes ### III.1 Dipolar oscillations The collisional state of the gas is revealed by the behavior of the frequencies and damping rates of collective excitations as functions of the collision frequency. Briefly, for dipolar modes the hydrodynamic behavior is signalled by a common frequency of oscillation of the two species and by a decrease of the damping rate with increasing collision frequency. Conversely the collisionless regime, which is attained at low collisionality, is characterized by different oscillation frequencies and by growing damping rates. The collision rate can be evaluated either from a numerical simulation which actually counts the number of collisions at each time step, or by direct integration of the collision integral over momenta. Pauli blocking strongly quenches collisions at the temperatures of present interest and its handling in numerical studies requires suitably adapted methods of Monte Carlo sampling Toschi et al. (2003); Capuzzi et al. (2004); Succi et al. (2004). In the numerical simulation we excite dipolar modes by initially shifting the fermionic density profile by $`a_{ho}=\sqrt{\mathrm{}/(m_F\omega _{F,})}`$ in either the axial or radial direction ($`\widehat{z}`$ and $`\widehat{x}`$, say) and then fit the time evolution of the fermionic center-of-mass coordinates with the functions $`\mathrm{cos}(\mathrm{\Omega }_it+\varphi )\mathrm{exp}(\gamma _it)`$. From these fits we extract the oscillation frequencies $`\mathrm{\Omega }_i`$ and the damping rates $`\gamma _i`$ along the two directions. The results obtained from the numerical solution of the Vlasov-Landau equations are shown in Fig. 1. The error bars in the plot have been estimated from the standard deviations of $`\mathrm{\Omega }_i`$ and $`\gamma _i`$ in the fitting process. Starting from $`\lambda =1`$ and increasing the anisotropy towards an elongated cigar-shaped trap, the frequencies of the dipolar oscillations show different behaviors along the axial and the radial direction. The value of $`\mathrm{\Omega }_{}/\omega _{F,}`$ increases slightly towards unity as $`\lambda `$ is decreased, whereas $`\mathrm{\Omega }_z/\omega _{F,z}`$ decreases with $`\lambda `$ towards an appreciably lower value. At the same time both damping rates tend to vanish, although that for axial oscillations appears to go through a broad maximum before doing so. These behaviors indicate that in strongly elongated traps the dynamics of the Fermi gas is collisionless in radial dipolar oscillations, but collisional in axial ones. The coexistence of collisionless and hydrodynamic behaviors in the dipolar oscillations as presented above is supported by the solution of scaling equations in the classical limit. Within a classical model one can write a set of equations for the coupled motions of the centers of mass of fermions and bosons in the small-oscillation regime Gensemer and Jin (2001); Ferlaino et al. (2003). If we include mean-field effects, the equations of motion in the $`i`$-direction for the center-of-mass coordinates and for the relative coordinates are $$\{\begin{array}{cc}\hfill \ddot{x}_{\mathrm{CM},i}& =\omega _{\text{hd},i}^2x_{\mathrm{CM},i}M_r\mathrm{\Delta }\omega ^2x_{r,i}/M\hfill \\ \hfill \ddot{x}_{r,i}& =\mathrm{\Delta }\omega ^2x_{\mathrm{CM},i}\omega _{r,i}^2x_{r,i}+\omega _{\mathrm{mf},i}^2x_{r,i}\omega _Q\dot{x}_{r,i}\hfill \end{array}$$ (5) where $`x_{\text{CM},i}=(m_FN_Fx_{F,i}+m_BN_Bx_{B,i})/M`$ and $`x_{r,i}=x_{F,i}x_{B,i}`$. We have defined $`\mathrm{\Delta }\omega ^2=\omega _{F,i}^2\omega _{B,i}^2`$ and $`\omega _{r,i}^2=(m_BN_B\omega _{F,i}^2+m_FN_F\omega _{B,i}^2)/M`$ with $`M=m_FN_F+m_BN_B`$ and $`M_r=m_Bm_FN_BN_F/M`$. The hydrodynamic frequencies are given by $$\omega _{\text{hd},i}=\left(\frac{m_FN_F\omega _{F,i}^2+m_BN_B\omega _{B,i}^2}{M}\right)^{1/2}$$ (6) and the collisional frequency is $`\omega _Q=4Q(m_B/N_F+m_F/N_B)/(3m_F+3m_B)`$, $`Q`$ being the total number of collisions per unit time defined as $$Q=\frac{\sigma }{\pi ^2}\frac{N_FN_B}{k_BT}\left(\frac{k_xk_yk_z}{m_r}\right)^{1/2}.$$ (7) In Eq. (7) $`k_i=k_{F,i}k_{B,i}/(k_{F,i}+k_{B,i})`$ are the effective oscillator constants, with $`k_{j,i}=m_j\omega _{j,i}^2`$. In the large-$`Q`$ limit Eqs. (5) predict that the relative motion of the two clouds is overdamped, so that fermions and bosons oscillate together at the hydrodynamic frequencies. For the system parameters that we are using these are $`\omega _{\text{hd},i}=0.9942\omega _{F,i}`$. While due to the low number of impurities the value of $`\omega _{\text{hd},i}`$ is very close to the bare trap frequency $`\omega _{F,i}`$, their difference can be amplified by, e.g., increasing the mass of the impurities or varying the trapping frequency. The mean-field correction to the bare frequency of the relative motion in Eq. (5) is $$\omega _{\mathrm{mf},i}^2=\frac{g}{M_r}d^3r\frac{n_0^{\left(F\right)}}{x_i}\frac{n_0^{\left(B\right)}}{x_i},$$ (8) where $`n_0^{(F,B)}(𝐫)`$ are the equilibrium density profiles. Equation (8) extends the result of Ref. Vichi and Stringari (1999) to a general gaseous mixture. According to Eqs. (5) the mean-field correction is more important in the collisionless limit and does not affect the value of the hydrodynamic frequency. However, it shifts the value of $`Q^{\text{lock}}`$ at which the two clouds become glued together. The locking point can be estimated by looking for an overdamped oscillation in the solution for the relative motion in Eqs. (5). If we neglect the coupling between the center-of-mass and the relative motions we find $$\sqrt{\omega _{r,i}^2\omega _{\mathrm{mf},i}^2}\frac{\omega _Q}{2}\|_{Q^{\text{lock}}}.$$ (9) We have estimated a value of $`\omega _{\text{mf},i}^2/\omega _{F,i}^28\times 10^3`$ for the isotropic trap and verified that this ratio is slightly decreasing with $`\lambda `$. Therefore the mean-field correction is negligible in the present case and the locking point is fixed by the bare frequency of the relative coordinate, namely $`\omega _Q\|_{Q^{\text{lock}}}2\omega _{r,i}1.4\omega _{F,i}`$. The classical model does not contain the effects of Pauli blocking at low temperature, so that we have taken $`Q`$ as a single fitting parameter at all values of $`\lambda `$. The results obtained from Eqs. (5) with $`Q=55\overline{\omega }_F`$ are reported in Fig. 1 and compared with the numerical results for the dependence of the mode frequencies and damping rates on the anisotropy parameter. At $`\lambda =0.05`$ the gas is already very close to the hydrodynamic regime in its motions along $`\widehat{z}`$, since the value of the oscillation frequency obtained in the simulation is close to $`\omega _{\text{hd},i}`$ in spite of the low concentration of impurities. The overall agreement between the simulation and the model is fairly good, and the deviations may be attributed to low-temperature effects and to shape deformations of the distributions not entering Eqs. (5). Finally, it is worthwhile stressing that even though the effects on the dipolar oscillations due to the scattering with impurities are small, they can be greatly increased by tuning the experimental parameters. For instance, if we consider the parameters of the recent experiment by Takasu et al. Takasu2003 and take 100 <sup>174</sup>Yb atoms as impurities, the difference between the oscillation frequencies along the axial and radial directions as shown in Fig. 1 will be of about $`75\%`$ in the limit of small $`\lambda `$. Therefore, we may expect that the effects here described will become easily observable in the near future when mixtures of <sup>40</sup>K and <sup>174</sup>Yb will be experimentally realized. ### III.2 Monopolar and quadrupolar oscillations Dipolar modes are a simple example of collective modes that can be experimentally analyzed and whose features expose the collisional state of the gas. As illustrated by the classical model, these modes can be viewed mainly as coupled oscillations of the centers of mass of the two clouds in the absence of appreciable deformations of their shapes. Other collective modes can be explored by deforming the clouds in different ways. The lowest-frequency modes of this type are the surface monopolar ($`l=0`$) and quadrupolar ($`l=2`$) modes, and these can be excited by small deformations of suitable symmetry. While in an isotropic trap the monopolar mode and the $`\mathrm{}_z=0`$ quadrupolar mode can be independently excited, in an axially symmetric trap the angular momentum is not a good quantum number and the two oscillations with $`\mathrm{}_z=0`$ are coupled to each other. As the collisionality of the fermion gas is increased, the transition to the hydrodynamic regime in strongly elongated traps manifests itself through changes in the $`\mathrm{}_z=0`$ mode frequencies from $`2\omega _{F,z}`$ to $`\mathrm{\Omega }_{\text{hd}}=\sqrt{12/5}\omega _{F,z}`$ and from $`2\omega _{F,}`$ to $`\sqrt{10/3}\omega _{F,}`$ Vichi (2000); Griffin et al. (1997); Guéry-Odelin et al. (1999). A similar transition can be expected in the present system, where the impurities act to transfer momentum between noninteracting fermions. We have monitored several dynamical averages of the fermion cloud in the course of the simulation and in particular we have analyzed the time evolution of $`\chi (t)=2v_{F,z}^2v_{F,}^2`$ after compression of the cloud density by about 5%. This average measures the anisotropy of the velocity distribution and is the responsible for the coupling between monopole and quadrupole excitations Guéry-Odelin et al. (1999). At variance from the dipole modes, in the monopole and quadrupole modes the radial and axial motions are strongly coupled and this requires that we simultaneously follow the dynamics of the gas on two quite different time scales. In Fig. 2 we plot the value of the lower $`\mathrm{}_z=0`$ mode frequency obtained from the Fourier transform of $`\chi (t)`$. At very large anisotropy this peak frequency suddenly drops from the collisionless value $`2\omega _{F,z}`$ towards the hydrodynamic value $`\sqrt{12/5}\omega _{F,z}`$, whereas the higher peak frequency (not shown in Fig. 2) stops at its collisionless value $`2\omega _{F,}`$. Again collisionless and hydrodynamic behaviors coexist in a strongly elongated trap. The behavior of monopolar and quadrupolar modes with finite relaxation time $`\tau _{mq}`$ in a gas of interacting fermions inside an anisotropic trap has been studied by Vichi Vichi (2000), who derived an implicit polynomial equation predicting a very steep downturn of the lower mode frequency in systems with $`\overline{\omega }_F\tau _{mq}1`$. The solid line in Fig. 2 shows the result of fitting Vichi’s model to our data with the choice $`\overline{\omega }_F\tau _{mq}=15`$. The sharp transition to the collisional regime is due to the fact that the impurities mediate the fermion-fermion scattering and lead to an effective relaxation time which is larger than that involved in the impurity-fermion scattering. Even though the effect is limited to large anisotropies, it could become observable at moderate anisotropies by increasing the numbers of particles or the strength of the boson-fermion repulsion. ## IV Expansion dynamics Many experiments extract information on the properties of an ultracold trapped gas after it been allowed to undergo free expansion. This improves the spatial resolution of in-situ measurements and also gives access to the momentum distribution of the atoms. A fermion gas in the collisionless regime becomes spatially isotropic as it expands, regardless of the initial anisotropy of its density profile, whereas in the hydrodynamic regime the density profile inverts its aspect ratio during expansion. Furthermore, since the number of collisions diminishes as the gas expands, the hydrodynamic picture will become invalid when there are not enough collisions to sustain local equilibrium. In the numerical simulation of the expansion, starting from the equilibrium density profiles of the trapped gaseous mixture we switch off the confinement and allow evolution according to the VLE. During the evolution we adaptively change the size of the computational domain in order to ensure that all particles are inside it. This permits us to study the free expansion of the cloud for long periods of time up to expansion ratios $`b_i(t)=R_i(t)/R_i(0)`$ of a few hundreds, with $`R_i(t)=x_{F,i}^2(t)^{1/2}`$ being the width of the cloud in the $`i`$-direction. In Fig. 3 we show the asymptotic value $`R_{}/R_z`$ of the aspect ratio of the fermion cloud as a function of the anisotropy $`\lambda `$. The aspect ratio in the case of isotropic or strongly anisotropic confinement is equal to unity as in the collisionless regime, while for intermediate values of $`\lambda `$ the scattering against the impurities makes $`R_{}/R_z`$ deviate from unity. This behavior can be seen as a consequence of the competition between the time scales associated to axial and radial motions, as explained below. We have evaluated the total numbers $`𝒬_{}`$ and $`𝒬_z`$ of collisions occurring during a lapse of time from $`t=0`$ to $`t_{}=1/\omega _{F,}`$ and to $`t_z=1/\omega _{F,z}`$, respectively. Assuming that the collision rate $`Q(t)`$ scales in time with the volume occupied by the gas and that the expansion dynamics is close to that of a collisionless system, $`𝒬_{}`$ and $`𝒬_z`$ can be written as $`𝒬_{,z}`$ $`=`$ $`{\displaystyle _0^{t_{},t_z}}𝑑tQ(t)`$ $`=`$ $`Q(0){\displaystyle _0^{t_{},t_z}}𝑑t(1+\omega _{F,}^2t^2)^1(1+\omega _{F,z}^2t^2)^{1/2}.`$ In Eq. (IV) we have taken $`Q(t)=Q(0)/(b_{}^2b_z)`$ and set $`b_i(t)=(1+\omega _{F,i}^2t^2)^{1/2}`$ Pedri et al. (2003). The integral in Eq. (IV) is straightforward and yields the results shown in Fig. 4. The following points should be noted: (i) the characteristic time $`t_{}`$ for the radial expansion and consequently $`𝒬_{}`$ decrease with increasing the anisotropy; and (ii) the integrand of Eq. (IV) for $`𝒬_z`$ becomes negligible for $`t_{}tt_z`$ as the fermion density drops due to rapid expansion in the radial direction, so that $`𝒬_z`$ increases at first with decreasing $`\lambda `$, reaches its maximum at $`\lambda 0.4`$, and then rapidly drops as the density of the expanding cloud rapidly vanishes. The non-monotonic behavior of $`R_{}/R_z`$ as a function of $`\lambda `$ in Fig. 3 can be understood from the features of $`𝒬_{}`$ and $`𝒬_z`$ in Fig. 4. The number $`𝒬_{}`$ of collisions diminishes with decreasing $`\lambda `$ while $`𝒬_z`$ increases, so that collisions play different roles in the axial and radial expansion. Both $`𝒬_{}`$ and $`𝒬_z`$ are dropping at large anisotropies, so that collisionless behavior is emerging in the ballistic expansion from a strongly elongated trap. The behavior of the expanding gas as a function of $`\lambda `$ can also be analyzed by means of scaling equations. In this approach Pedri et al. (2003) the expansion of a fermion gas with finite collisionality can be described by the equations $$\{\begin{array}{cc}\hfill \ddot{b}_i\omega _{F,i}^2\frac{\theta _i}{b_i}& =0\hfill \\ \hfill \dot{\theta }_i+2\frac{\dot{b}_i}{b_i}\theta _i^2& =\frac{1}{\tau (\theta _i,b_i)}\left(\theta _i\frac{1}{3}\underset{j}{}\theta _j\right)\hfill \end{array}$$ (11) where $`\theta _i(t)`$ are effective temperatures along the two spatial directions, $`\tau (\theta _i,b_i)=\tau _0(_jb_j)(\frac{1}{3}_k\theta _k)^{1/2}`$ and $`\tau _0`$ is the collision time at $`t=0`$. The collisionless and hydrodynamic limits correspond to taking $`\tau =\mathrm{}`$ and $`\tau =0`$ in Eqs. (11), respectively. We have taken the collision time $`\tau _0`$ in Eqs. (11) as a fitting parameter to represent the indirect fermion-fermion scattering induced by the impurities. The solution of Eqs. (11) with the choice $`\tau _0=6/\overline{\omega }_F`$ reproduces the non-monotonic behavior of the aspect ratio as shown in Fig. 3. Finally, we have checked that the qualitative behavior shown in Fig. 4 for $`𝒬_{}`$ and $`𝒬_z`$ does not depend on the assumption of a collisionless evolution of the gas and persists on describing its expansion with any value of $`\tau _0`$. The magnitude of the deviation from unity of the aspect ratio of the expanding fermions depends on the value of $`\tau _0`$ and can be increased by changing the fermion-impurity interaction via Feshbach resonances Inouye2004a or by performing an out-of-equilibrium experiment as described in Ref. Capuzzi et al. (2004). ## V Summary and concluding remarks We have studied the low-lying surface modes and the ballistic expansion of a spin-polarized Fermi gas interacting with thermal impurities as functions of the anisotropy of its confinement in a cigar-shaped harmonic trap. We have solved for this purpose the Vlasov-Landau equations for the dynamics of the mixture and compared the results with simple scaling equations containing collision-time fitting parameters. The results show that for large anisotropies a collisionless behavior in the radial dipolar oscillations and a hydrodynamic behavior in the axial ones are simultaneously established. For monopolar and quadrupolar excitations we have observed a collisionless spectrum irrespectively of the strength of the anisotropy parameter, except for extremely large anisotropies where the frequency of the lower $`\mathrm{}_z=0`$ mode decreases towards the hydrodynamic value. On the other hand, during ballistic expansion the two different time scales for radial and axial motions enter into competition. The result is that the expansion of a strongly anisotropic cloud is essentially collisionless due to the rapid drop of the particle density, whereas at intermediate values of the anisotropy the aspect ratio is sensitive to the collisions. The analysis presented here has only concerned slightly doped fermions in a collisionless-to-intermediate scattering regime as, e.g., that attained in the experiments at LENS. The effects of the trap anisotropy could be further enhanced exploiting the rich variety of experimental set-ups and trapped isotopes available. It would be also interesting to extend our work to strong-interaction situations among fermions and impurities or among fermions in two-component Fermi gases approaching the unitary limit. ###### Acknowledgements. This work has been partially supported by an Advanced Research Initiative of Scuola Normale Superiore di Pisa and by the Istituto Nazionale di Fisica della Materia within the Initiative “Calcolo Parallelo”.
warning/0506/cond-mat0506292.html	ar5iv	text	# Coherence Network in the Quantum Hall Bilayer ## Abstract Recent experiments on quantum Hall bilayers near total filling factor 1 have demonstrated that they support an “imperfect” two-dimensional superfluidity, in which there is nearly dissipationless transport at non-vanishing temperature observed both in counterflow resistance and interlayer tunneling. We argue that this behavior may be understood in terms of a coherence network induced in the bilayer by disorder, in which an incompressible, coherent state exists in narrow regions separating puddles of dense vortex-antivortex pairs. A renormalization group analysis shows that it is appropriate to describe the system as a vortex liquid. We demonstrate that the dynamics of the nodes of the network leads to a power law temperature dependence of the tunneling resistance, whereas thermally activated hops of vortices across the links control the counterflow resistance. Introduction– Over the past decade, there has been accumulating evidence that excitonic superfluid-like keldysh properties may be present in quantum Hall bilayer systems bilayers when the filling factor of each layer is near $`\nu =1/2`$ and the layers are sufficiently close. The relevance of excitons to this system may be understood by using a filled Landau level in a single layer as a starting reference state; equal densities in the layers may then be reached by creating a particle-hole condensate fertig89 ; eisenstein . It has long been suspected that such a state might have superfluid properties wen . The most dramatic superfluid-like properties of this system are observed when the layers are separately contacted and a current is passed from one layer to the other. For contacts at opposite ends of the sample, the resulting tunneling conductance $`\sigma _T`$ has a sharp resonance precisely at zero voltage spielman1 , which is reminiscent of the DC Josephson effect. When the contacts are on the same side of the sample and the layers on the opposite side are short-circuited, so that currents in the two layers flow in opposite directions \[“counterflow” (CF)\], the voltage drop in either layer along the direction of flow tends to zero in the low temperature limit kellogg ; tutuc . This may be understood as nearly dissipationless flow of electron-hole pairs. Interestingly, the temperature dependences of these two behaviors are very different: the tunneling resistance $`1/\sigma _T`$ falls much more slowly than the dissipative in-plane resistance. The “imperfect” superfluidity apparent in these experiments suggests that a good starting point for understanding this system is in terms of a condensed state with interlayer coherence. However, a true two-dimensional superfluid should have a power law $`IV`$ in CF below some critical Kosterlitz-Thouless temperature $`T_{KT}`$, whereas an Ohmic response is observed at all available temperatures. Moreover, a true DC Josephson effect would yield an infinite tunneling conductance below $`T_{KT}`$ rather than just a sharp resonance. It is also remarkable that in both the tunneling and CF geometries, current appears to flow over the entire length of the system, whereas for a true superfluid the current would decay within a Josephson length of the edge prange , which is far smaller. The linear response of the quantum Hall bilayer, exhibiting dissipationless transport, if at all, only at zero temperature, is a new transport regime for condensed matter systems. Most theoretical studies of this system suggest that disorder must be crucially involved in the mechanism(s) that lead to this behavior balents ; hafjp2 ; wang ; sheng . In this work, we describe a state where the experimentally observed behavior emerges naturally, and which is indeed induced by the disorder environment of the quantum Hall bilayer. We will argue below that the system breaks up into large regions in which the interlayer coherence is small or absent, separated by narrow regions in which the coherence is relatively strong. The resulting structure is a coherence network, with long, quasi-one-dimensional links connected at nodes. We will demonstrate that dissipation in tunneling may be understood in terms of the phase dynamics of the nodes, whereas dissipation in CF occurs due to thermally activated hops of vortices across the links. Further experimental consequences of our model are discussed below. Nonlinear Screening and the Coherence Network – We start with the Efros model for nonlinear screening of potential fluctuations in a quantum Hall system efros . This screening generates puddles of quasiholes or quasielectrons, which are locally dense and mobile enough to destroy the incompressibility of the quantum Hall fluid. The puddles are separated by narrow ($`\mathrm{}_0`$, with $`\mathrm{}_0=\sqrt{\mathrm{}c/eB}`$ the magnetic length) incompressible strips which follow the equipotentials of the total effective potential, and within which there is a charge gap. The strips form links which join together at nodes – associated with saddle points of the disorder potential – to form a random network. For the bilayer system, the incompressible state at total filling factor $`\nu =1`$ has a phase degree of freedom $`\theta `$ that may be understood as the wavefunction phase for excitons in the incompressible regions of the system. The conjugate variable to this phase angle is the interlayer charge imbalance, which may be denoted as $`S_z`$. For equal layer populations, on average $`<S_z>=0`$, and the system becomes analogous to an easy-plane Heisenberg ferromagnet moon . Vortices of the phase field (“merons”) carry both an electric charge ($`\pm e/2`$ for either vorticity) and an interlayer electric dipole moment moon ; lee ; huse . Thus, in this context the links and nodes of the network carry a phase degree of freedom, while the puddles are flooded with vortex-antivortex pairs puddles . The presence of meron-antimeron pairs at link edges implies that the phase angle $`\theta (𝐫)`$ turns over many times along a given link. When tunneling is included, these overturns bear a close resemblance to kink solitons of the sine-Gordon model moon . As discussed below, the phase dynamics of the links is conveniently described in terms of these solitons. Our picture of this coherence network is illustrated in Fig. 1. Renormalization Group Analysis – For counterflow experiments, dissipation can only occur when merons – i.e., vortices – are mobile and unbound huse . For this to occur, the vortices in the coherence network must be in a liquid phase; i.e., their discreteness should not inhibit their ability to screen “charge” (i.e., enforced vorticity) at long distances. We begin with a classical model on a square lattice network, $`H_{CN}={\displaystyle \underset{𝐫}{}}\left\{{\displaystyle \frac{1}{2}}K{\displaystyle \underset{\mu =x,y}{}}\left[2\pi m_\mu (𝐫)_\mu \theta (𝐫)+A_\mu (𝐫)\right]^2h\mathrm{cos}\theta (𝐫)\right\},`$ where $`m_{x,y}(𝐫)`$ are integer degrees of freedom representing the number of solitons on the link extending from the node at $`𝐫`$ in the $`x`$ or $`y`$ direction, and $`A_\mu (𝐫)`$ is a Gaussian random variable obeying $`<A_\mu (𝐫)A_\nu (𝐫^{})>=\mathrm{\Delta }\delta _{\mu ,\nu }\delta _{𝐫,𝐫^{}}`$ which models the effect of disorder. The parameter $`h`$ is proportional to the interlayer tunneling matrix element, which explicitly breaks the $`U(1)`$ symmetry in energetically favoring $`\theta =0`$. It is convenient to recast the problem explicitly in terms of vortices, which can be done by standard techniques jose , leading to $`Z=𝒟\varphi _{\{B\}}e^{H_V}`$, with $`H_V=2\pi ^2K{\displaystyle \underset{𝐫}{}}\{\|\varphi (𝐫)+\widehat{x}B(𝐫)+𝐀(𝐫)\|^2`$ $`y_h\mathrm{cos}2\pi \varphi (𝐫)+E_c\left({\displaystyle \frac{B(𝐫)}{y}}\right)^2\}.`$ (1) In this representation $`B`$ is an integer field that resides on vertical bonds, and the vorticity is $`m(𝐫)=B/y`$. Finally, $`E_c`$ represents a core energy for these vortices, and $`y_h=e^{k_BT/h}`$. This form of $`H_V`$ is appropriate when $`y_h<<1`$, which is easily satisfied under experimental conditions. For $`y_h=0`$ and $`𝐀=0`$, $`\varphi `$ may be integrated out to recover the usual Coulomb gas. To facilitate the renormalization group (RG) analysis, we replace the integer field $`B(𝐫)`$ with a continuous field $`b(𝐫)`$ and add a term $`y_b_𝐫\mathrm{cos}[2\pi b(𝐫)]`$ that tends to keep $`b`$ at its intended integer values. By matching low energy configurations for the integer and continuous field versions of the theory, one may show $`y_bE_c`$ for small $`E_c`$ and $`y_h`$. The crucial point is that the irrelevance of $`y_bE_c`$ will indicate a liquid state of vortices. In performing the RG, $`b`$ is rescaled according to $`b^{}(𝐫^{})e^{\mathrm{}}=b(𝐫)`$, where $`e^{\mathrm{}}`$ is the rescaling factor, in order to keep the Gaussian part of the Hamiltonian of fixed form. As $`b`$ shrinks with rescaling, it is natural to expand $`\mathrm{cos}[2\pi b(𝐫)]`$ in a power series in $`b`$, and treat the fourth and higher order terms perturbatively. The quadratic term, however, introduces a contribution to the fixed point of the form $`\frac{1}{2}r_𝐫b(𝐫)^2`$. A key observation is that this term limits large separations for vortices, so that if $`r0`$ at the fixed point, the system is in a bound vortex state deconfine1 ; deconfine2 . In the clean limit ($`𝐀=0`$), it has been demonstrated that this happens if the temperature is too low, $`E_c`$ is too large, or if $`h`$ is too large. By contrast, for $`r=0`$, states of arbitrarily large “vortex dipole moment” may be found in the thermodynamic limit, so that the vortices are effectively unbound deconfine2 . To generalize this to the disordered case, we use the replica trick before proceeding with the RG analysis. In the $`n0`$ limit, the flow equations become, to lowest order in $`E_c`$ and $`y_h`$, $`{\displaystyle \frac{d}{d\mathrm{}}}\left({\displaystyle \frac{r}{K}}\right)`$ $`=`$ $`e^2\mathrm{}\sqrt{{\displaystyle \frac{r/K}{1+\xi ^2\mathrm{\Lambda }^2}}}\mathrm{\Lambda }^2\left[{\displaystyle \frac{1}{K}}+2\pi ^2\mathrm{\Delta }\right]`$ $`{\displaystyle \frac{d\xi ^2}{d\mathrm{}}}=`$ $`2\xi ^2`$ $`{\displaystyle \frac{dy_h}{d\mathrm{}}}=(2\eta /\sqrt{Kr})y_h`$ where $`\mathrm{\Delta }`$ is the disorder strength, $`\mathrm{\Lambda }=\pi /a_0`$ is a wavevector cutoff, $`a_0`$ is the distance between nodes, $`\xi `$ is a length scale with initial value $`\xi (\mathrm{}=0)=\sqrt{E_c/4\pi ^2K}a_0`$, and the initial value of $`r(\mathrm{}=0)E_c`$. The parameter $`\eta `$ is a number of order unity. To this order in perturbation theory, the parameter $`K`$ remains constant. For large disorder $`\mathrm{\Delta }`$ and/or small $`K`$ sheng ; stiffness , it is apparent that $`r,y_h0`$. Since all signs of the discreteness of the underlying the vortex density has scaled away at this fixed point, it is a liquid. In this phase neither the stiffness $`K`$ nor the symmetry-breaking field $`h`$ are effective at preventing phase overturns along a column of nodes (say, at the edge of the system) if their phase angles are driven by an arbitrarily small torque wei . Moreover, the irrelevance of $`y_h`$ means the tunneling term $`h\mathrm{cos}\theta `$ may be treated perturbatively in the vortex liquid phase. We now exploit this observation to understand the transport properties of the coherence network. Dissipation in Transport: Tunneling and Counterflow – We can understand transport in the coherence network based on the following considerations. (i) The total current in the sample is the sum of a CF current and the quantum Hall edge current. Dissipation in the latter is negligible in both the tunneling and CF geometries we consider here, and so is ignored in what follows. (ii) We assume the solitons are unpinned in the links, as is appropriated for slowly varying disorder, and treat them as an elastic system. The displacement field $`u(x)`$ is related to the counterflow current via $`I_{CF}=\rho _s_x\theta \stackrel{~}{K}_xu`$ (with $`\stackrel{~}{K}\rho _s`$). From the Josephson relation, it follows that the interlayer electric potential is $`V_{int}=\frac{\mathrm{}}{e}_t\theta =\frac{h}{e}_tu/b_0`$, where $`b_0`$ is the average spacing of the solitons. A simple effective Gaussian theory may be written for the solitons on the links unpub , whose form would control dissipation at sufficiently high frequencies. Such a purely quadratic theory however cannot dissipate energy at zero frequency maslov-stone . (iii) CF current can be degraded by hopping of merons across the link. This is a thermally activated process with energy $`\mathrm{\Delta }_{link}`$, because the merons must cross the barrier of the incompressible link. (iv) The nodes are relatively meron-free regions with a high stiffness, and can be treated as rigid rotors which are subject to torques $`F_{link}I_{CF}^{link}`$ from each attached link, and a force $`h\mathrm{sin}\theta `$ due to the tunneling. Consider first the tunneling geometry. Since the current flows into (say) the top layer on the left and leaves via the bottom layer on the right, the CF current points in opposite directions at the two ends of the sample. From (ii) above this implies the currents cause the underlying phase angles in the links and nodes to rotate in the same direction, and may be modeled by forcing in solitons at one sample end and removing them at the other, as illustrated in Fig. 2. The dynamics of a typical node may be described by a Langevin equation $$\mathrm{\Gamma }\frac{d^2\theta }{dt^2}=\underset{links}{}F_{link}\gamma _0\frac{d\theta }{dt}h\mathrm{sin}\theta +\xi (t).$$ (2) Here $`\mathrm{\Gamma }`$ is the effective moment of inertia of a rotor, proportional to the capacitance of the node, $`\xi `$ is a random (thermal) force, and $`\gamma _0`$ is the viscosity due to dissipation from the other node rotors in the system. For a small driving force, the node responds viscously, and the resulting rotation rate determines the rate of flow of solitons via $`\gamma \dot{\theta }=F_{link}`$. From (ii) and (iv) above one sees that the viscosity $`\gamma `$ is proportional to the tunneling conductance $`\sigma _T`$ of the system. For $`k_BTh`$ one may show the viscosity for an individual node to be dieterich $$\gamma =\gamma _0+\mathrm{\Delta }\gamma =\gamma _0+\sqrt{\frac{\pi }{2\mathrm{\Gamma }}}\frac{h^2}{(k_BT)^{3/2}}.$$ (3) As each node contributes the same amount to the total viscosity, the total response of the system to the injected CF current obeys $$I_{CF}N_{nodes}\mathrm{\Delta }\gamma \frac{e^2V_{int}}{\mathrm{}}=\sigma _TV_{int}$$ (4) Note that because the nodes respond viscously, the tunneling conductance is proportional to the area of the bilayer. This is so far consistent with experiment but difficult to understand without the network geometry posited here balents . The dependence of the tunneling conductance on temperature, system area, and $`h`$ each constitute an experimentally falsifiable prediction of the coherence network model of the quantum Hall bilayer com1 ; com2 . Now consider the CF geometry. This is created by short-circuiting the layers at one end of the sample, which we will take to be on the right. There the rotor angles should be subject to free boundary conditions, leading to $`F\dot{\theta }\dot{u}=0`$. The injected CF current on the left leads to a compression of the solitons (or a rarefaction of antisolitons). In a uniform superfluid, such a force would be fully balanced by the effective force due to tunneling, and the CF current would flow without dissipation. In the coherence network, however, at any $`T0`$, merons are able to thermally hop across the links to create or destroy solitons, and the other (unpinned) solitons move in response, leading to dissipation. In equilibrium, the rates of creation and annihilation of solitons in a link obey detailed balance. However, the presence of the “force” $`FI_{CF}`$ destroys this balance, and leads to a net meron current $$I_{meron}=\frac{\zeta }{e}I_{CF}e^{\mathrm{\Delta }_{link}/k_BT},$$ (5) where $`\zeta `$ is a dimensionless number depending on the details of the link disorder. If, for example, there is a net annihilation of solitons with time on a given link, solitons will move in from the left to take their place. The rate of change of $`\theta `$ due to such processes at a point inside the sample a distance $`x`$ from the left end is $$\frac{d\theta (x)}{d\tau }=2\pi \underset{x}{\overset{L}{}}I_{meron}(x^{})𝑑x^{}.$$ (6) For temperatures satisfying $`hk_BT\mathrm{\Delta }_{link}`$, the degrading of $`I_{CF}`$ due to the nodes may be ignored. In this case, Eqs. 5 and 6 together with the Josephson relation yield an activated longitudinal resistance measured between probes $`\mathrm{\Delta }x`$ apart on a single layer of the form $$R_{xx}=(\mathrm{\Delta }x)\zeta \frac{h}{e^2}e^{\mathrm{\Delta }_{link}/k_BT}$$ (7) which is consistent with experiment. In the CF geometry, the meron current $`I_{meron}`$ also creates a Hall voltage within individual layers because merons carry an interlayer dipole moment, leading to nonequilibrium charge imbalances at the upper and lower edges of the sample of opposite sign. Assuming a $`T`$-independent resistive relaxation at the edges (where the electrons are likely to be compressible), this creates a Hall voltage of opposite sign on the two layers which is activated, with the same activation energy as the longitudinal resistance. This is consistent with measurements in electron samples kellogg ; holes . Experimental Consequences – As shown in this work, the imperfect superfluidity observed in quantum Hall bilayers may be understood if we assume the electrons have organized into a coherence network. Several experimental tests of this hypothesis are possible. These include the temperature and size dependence of the tunneling conductance discussed above. Another interesting test would be measurement of the finite frequency tunneling conductance, for which the response of the links should become accessible. An additional possibility is the introduction of an artificial weak link across the direction of counterflow current. This would create a favored channel for soliton destruction/creation that should lead to a decrease in the activation energy for the CF resistance and an enhanced Hall voltage at its endpoints. The authors would like to thank S. Das Sarma, J.P. Eisenstein, M. Shayegan, and E. Tutuc for useful discussions. This work was supported by the NSF via Grant Nos. DMR0454699 (HAF) and DMR0311761 (GM).
warning/0506/gr-qc0506039.html	ar5iv	text	# The robustness of a many-body decoherence formula of Kay under changes in graininess and shape of the bodies ## 1 Introduction and Results In , one of us argued that several of the puzzles inherent in our current understanding of quantum and quantum-gravitational physics would appear to find a natural resolution if one were to postulate that quantum gravity is a quantum theory of a conventional type with a total Hilbert space $`H_{\text{total}}`$ which arises as a tensor product $$H_{\text{total}}=H_{\text{matter}}H_{\text{gravity}}$$ of a matter and a gravity Hilbert space, and with a total time-evolution which is unitary, but that, while, as would usually be assumed in a standard quantum theory, one still assumes that there is an ever-pure time-evolving “underlying” state, modelled by a density operator of form $$\rho _{\text{total}}=\|\mathrm{\Psi }\mathrm{\Psi }\|,$$ $`\mathrm{\Psi }H_{\text{total}}`$, at each “instant of time”, one adds the new assumption that the physically relevant density operator is not this underlying density operator, but rather its partial trace, $`\rho _{\text{matter}}`$, over $`H_{\text{gravity}}`$. With these assumptions, an initial underlying state of a closed quantum gravitational system with a low degree of matter-gravity entanglement would be expected to become more and more entangled as time increases – equivalently $`\rho _{\text{matter}}`$ will be subject to an ever-increasing amount of decoherence. In fact the von-Neumann entropy of $`\rho _{\text{matter}}`$ would be expected to increase monotonically, thus (when the theory is applied to a model for the universe as a whole) offering an explanation for the Second Law of Thermodynamics and (when the theory is applied to a model closed system consisting of a black hole sitting in an otherwise empty universe) offering a resolution to the Information Loss Puzzle. (It offers both of these things in that, by defining the physical entropy of $`\rho _{\text{total}}`$ to be the von-Neumann entropy of $`\rho _{\text{matter}}`$, one reconciles an underlying unitary time-evolution on $`H_{\text{total}}`$ with an entropy which varies/increases in time.) A second paper, , by the same one of us investigated the implications of the theory proposed in for the decoherence of ordinary matter in the non-relativistic, weak-gravitational-field regime. In particular, addressed the question of the appropriate description of the state, at some instant of time, of the center-of-mass degrees of freedom of a system of $`N`$ bodies, which, for simplicity, were taken to be identical. In ordinary quantum mechanics, such a state would be described by a many-body Schrödinger wave function $`\psi (𝐱^\mathrm{𝟏},\mathrm{},𝐱^𝐍)`$ in the usual matter Hilbert space $`H_{\text{matter}}`$, consisting of the (appropriately symmetrized, according to whether the bodies are treated as fermions or bosons or neither) $`N`$-fold tensor product of $`L^2(R^3)`$. argues that such a $`\psi `$ needs to be replaced by a total (entangled) matter-gravity state vector $`\|\mathrm{\Psi }`$ in a total matter-gravity Hilbert space which is the tensor product of $`H_{\text{matter}}`$ (defined as above) with a suitable Hilbert space $`H_{\text{gravity}}`$ for the modes of the quantized linearized gravitational field. Specifically, thinking of this tensor-product Hilbert space as the set of (appropriately symmetrized) square-integrable functions from $`R^{3N}`$ to $`H_{\text{gravity}}`$, $`\|\mathrm{\Psi }`$ is taken, in , to be the function $$(𝐱_1,\mathrm{},𝐱_N)\psi (𝐱^\mathrm{𝟏},\mathrm{}𝐱^𝐍)\|\gamma (𝐱^\mathrm{𝟏},\mathrm{}𝐱^𝐍)$$ where $`\|\gamma (𝐱^\mathrm{𝟏},\mathrm{}𝐱^𝐍)`$ is a certain (non-radiative) quantum coherent state (introduced in ) of the linearised gravitational field describing the Newtonian gravitational field due to the simultaneous presence of a body centred on each of the locations $`𝐱^\mathrm{𝟏},\mathrm{}𝐱^𝐍`$. (We remark that, in , for the purposes of calculating these coherent states, the bodies are modelled as classical mass-distributions.) In the non-relativistic, weak-gravitational-field regime, the projector $`\rho _{\text{total}}=\|\mathrm{\Psi }\mathrm{\Psi }\|`$ onto this $`\|\mathrm{\Psi }`$, is thus taken to be, to a very good approximation, the correct description of the “underlying” state of quantum gravity describing our many-body system in the sense described in the previous paragraph. The physically relevant density operator is therefore given, in this approximation, by the partial trace $`\rho _{\text{matter}}`$ of this $`\rho _{\text{total}}`$ over $`H_{\text{gravity}}`$ which, in position space is clearly given by $$\rho _{\text{matter}}(𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{})$$ $$=\psi (𝐱_1,\mathrm{},𝐱_N)\psi ^{}(𝐱_1^{},\mathrm{},𝐱_N^{})(𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{})$$ where the multiplicative factor $`(𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{})`$ is defined by $$(𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{})=\gamma (𝐱_1,\mathrm{},𝐱_N)\|\gamma (𝐱_1^{},\mathrm{},𝐱_N^{}).$$ (In , $``$ was written $`e^D`$ where $`D`$ was called the decoherence exponent.) Moreover, it was shown in that, in the case the bodies are taken to be balls with constant mass density and (we work in Planck units where $`G=c=\mathrm{}=1`$) total mass $`M`$ and radius $`R`$ (i.e. in the case the classical mass distributions representing the bodies are taken to be such balls) and in the case all the positions $`\{𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{}\}`$ of their centres of mass are much further away from one another than $`R`$ (we shall call this the well-spaced regime)<sup>1</sup><sup>1</sup>1In many applications, one expects the region of $`(\text{configuration space})\times (\text{configuration space})`$ where this condition doesn’t hold to be so small in comparison to size of the region where $`\psi \psi ^{}`$ is significantly big that no significant correction would be needed to results calculated on the assumption that the multiplicative factor is well-approximated by the $`_a`$ in all of $`(\text{configuration space})\times (\text{configuration space})`$., the multiplicative factor is well approximated by $`_a`$ where $`_a`$ is given by the explicit formula $$_a(𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{})=\underset{I=1}{\overset{N}{}}\underset{J=1}{\overset{N}{}}\left(\frac{\|𝐱_I^{}𝐱_J\|\|𝐱_I𝐱_J^{}\|}{\|𝐱_I𝐱_J\|\|𝐱_I^{}𝐱_J^{}\|}\right)^{12M^2}$$ (1) where it is to be understood that, in the cases $`I=J`$, the terms in the denominator $`\|𝐱_I𝐱_J\|\|𝐱_I^{}𝐱_J^{}\|`$ are to be replaced by $`R^2`$. To summarize, taking into account gravitational effects and then tracing over the gravitational field, as the general theory of prescribes that we should do, has, in the non-relativistic weak-gravitational-field regime, and on the assumption that the bodies are uniform mass balls, an overall effect which is equivalent to multiplying the usual quantum mechanical (pure) position-space density matrix $`\psi (𝐱_1,\mathrm{},𝐱_N)\psi ^{}(𝐱_1^{},\mathrm{},𝐱_N^{})`$ by a multiplicative factor, $``$, which, in the well-spaced regime, is well-approximated by $`_a`$ given by the formula (1) and it is the product of this pure density matrix with the multiplicative factor which is to be regarded as the physically relevant density operator. We remark that the formula (1) can also be written $$_a(𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{})=\left(\underset{K}{}\left(\frac{\|𝐱_K𝐱_K^{}\|}{R}\right)\underset{I<J}{}\left(\frac{\|𝐱_I^{}𝐱_J\|\|𝐱_I𝐱_J^{}\|}{\|𝐱_I𝐱_J\|\|𝐱_I^{}𝐱_J^{}\|}\right)\right)^{24M^2}$$ where the first product is taken over all $`K`$ from $`1`$ to $`N`$ and the second product is taken over all $`I`$ and $`J`$ from 1 to N, which satisfy $`I<J`$. We also remark that, in the case of a wave function for a single ball state, this prescription amounts to multiplying the density operator $`\psi (𝐱)\psi (𝐱^{})`$ by the multiplicative factor $$_a(𝐱;𝐱^{})=(\|𝐱𝐱^{}\|/R)^{24M^2}.$$ The problems we wish to address in the present paper are that, in the derivation, , of (1), the bodies are assumed to be constant mass-density balls whereas in applications, the bodies one is interested in will actually typically be “grainy”, and also not necessarily spherical. In fact, in one application (cf. the discussion of the Schrödinger-cat-like states in ) of the above formulae, the balls are interpreted as macroscopic balls of ordinary matter and real ordinary matter is of course grainy in that it is made out of atoms etc. and we may also be interested in lumps of ordinary matter with shapes other than balls. In another application, the above formulae (or rather their obvious generalization to states of many bodies where the masses and radii are not all equal) are interpreted as telling us how the standard non-relativistic quantum mechanics of closed systems of large numbers of atomic nuclei and electrons etc. gets modified according to the theory developed in and . In this latter application, the balls are taken to be models of nuclei and electrons etc. and again, of course, real nuclei and electrons are not actually uniform density balls, but will have their own grainy substructure and also need not be spherical. However one might hope that if constant-density balls were replaced by grainy, and possibly non-spherical, (classical) “lumps”, then the formula $`(\text{1})`$ for the multiplicative factor might nevertheless remain approximately correct provided we replace $`R`$ in $`(\text{1})`$ by a suitable effective radius $`R_{\text{eff}}`$. One might further hope that, in the case of a grainy ball, $`R_{\text{eff}}`$ would turn out to be of the same order of magnitude as $`R`$ and in the case of a grainy lump of some other shape, to be of the same order of magnitude as some measure of the lump’s typical linear size. (We shall refer below to the lump’s diameter without attempting to give a precise general definition to this notion.) If these hopes were fulfilled, then one might say that the formula $`(\text{1})`$ for $`_a`$ is robust under (classical) changes in graininess and in shape. We remark that we are continuing to assume here that our grainy, non-spherical, situations can, as far as their gravitational effects are concerned, still be modelled within the formalism of , as classical (albeit now no longer uniform and no longer spherical) mass distributions and we admit that, in principle, we should of course presumably work within an extension of the formalism of which allows for truly quantum (and also relativistic) descriptions of graininess (and non-sphericity). (Especially, such a quantum relativistic description may well be important in the second of our applications, to the grainy substructure of the proton etc.) While we shall not attempt to explore such quantum notions of graininess and non-sphericity in this paper, if robustness holds in the above sense under classical changes in graininess and in shape, it might also seem reasonable to guess that our formulae will remain robust even if graininess and changes in shape were taken into account in such a correct quantum way. In any case, the purpose of the present paper is to investigate whether and/or under what circumstances, these hopes are fulfilled in the case of one specific class of classical grainy models. Namely, models in which every one of the balls of mass $`M`$ and radius $`R`$ in formula $`(\text{1})`$ is replaced by a collection of $`n`$ little balls, each of mass $`m`$ and radius $`r`$ fixed at definite positions inside a lump-shaped region (i.e. some smooth closed bounded region of $`𝐑^3`$) of diameter approximately equal to $`2R`$ so that its centre of mass is located where the center (i.e. one of the positions $`𝐱_1,\mathrm{},𝐱_N`$) of the ball it replaces was located. We shall assume that the configurations of the balls within each of the lumps are Euclidean-congruent to one-another. However we shall not restrict the relative orientations of the different lumps; they are allowed to be arbitrarily rotated with respect to one another. We remark that obviously we have $`nm=M`$. Here we have in mind choices of values for the various parameters such that $$r\text{ any inter-small-ball spacing}andR\text{ any inter-lump spacing},$$ (2) where, by “any inter-lump spacing” we mean any of the distances between pairs of distinct elements of the set $`\{𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{}\}`$ and (assuming from now on that a definite system has been adopted for enumerating the little balls – as $`x_{I1},\mathrm{},x_{In}`$ etc. – within the $`I`$th lump etc.) by “any inter-small-ball spacing” we mean any of the distances between pairs of distinct elements of the set $`\{𝐱_{I1},\mathrm{},𝐱_{In};𝐱_{I1}^{},\mathrm{},𝐱_{In}^{}\}`$ for some/any $`I`$. Then (in all but an unimportantly small volume of $`(\text{configuration space})\times (\text{configuration space})`$ – cf. footnote above) formula $`(\text{1})`$ will clearly get replaced by $$_a^{\text{grainy}}(𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{})=\underset{I=1}{\overset{N}{}}\underset{i=1}{\overset{n}{}}\underset{J=1}{\overset{N}{}}\underset{j=1}{\overset{n}{}}\left(\frac{\|𝐱_{Ii}^{}𝐱_{Jj}\|\|𝐱_{Ii}𝐱_{Jj}^{}\|}{\|𝐱_{Ii}𝐱_{Jj}\|\|𝐱_{Ii}^{}𝐱_{Jj}^{}\|}\right)^{12m^2}$$ (3) where now it is to be understood that in the cases where $`I=J`$ and $`i=j`$, the terms in the denominator $`\|𝐱_{Ii}𝐱_{Jj}\|\|𝐱_{Ii}^{}𝐱_{Jj}^{}\|`$ are to be replaced by $`r^2`$. In line with $`(\text{2})`$, if we make the replacements: $$𝐱_{Ii}^{}𝐱_{Jj}𝐱_I^{}𝐱_J,\text{and}𝐱_{Ii}𝐱_{Jj}𝐱_I𝐱_J\text{except}\text{ when }I=J$$ (and similarly with primed and unprimed quantities interchanged) in $`(\text{3})`$, then we clearly expect to get a good approximation<sup>2</sup><sup>2</sup>2In fact this will clearly be true in the sense that if we scale the positions of the centres of mass of the lumps while not scaling the lumps themselves, then $`_a^{\text{grainy}}/_a`$ tends to $`1`$ as the scale tends to infinity. to $`_a^{\text{grainy}}`$. Making these replacements, we immediately get $`_a^{\text{grainy}}(𝐱_1,\mathrm{},𝐱_N;𝐱_1^{},\mathrm{},𝐱_N^{})`$ $`\left(\left({\displaystyle \frac{\underset{I}{}\underset{J}{}\|𝐱_I^{}𝐱_J\|^{n^2}\|𝐱_I𝐱_J^{}\|^{n^2}}{_{IJ}\|𝐱_I𝐱_J\|^{n^2}\|𝐱_I^{}𝐱_J^{}\|^{n^2}}}\right){\displaystyle \underset{i,j}{}}\left({\displaystyle \frac{1}{\|𝐱_i𝐱_j\|}}\right)^{2N}\right)^{12m^2}`$ $`=\left(\left({\displaystyle \frac{\underset{I}{}\underset{J}{}\|𝐱_I^{}𝐱_J\|\|𝐱_I𝐱_J^{}\|}{_{IJ}\|𝐱_I𝐱_J\|\|𝐱_I^{}𝐱_J^{}\|}}\right){\displaystyle \underset{i,j}{}}\left({\displaystyle \frac{1}{\|𝐱_i𝐱_j\|}}\right)^{2N/n^2}\right)^{12M^2}`$ where the final product involves all the positions $`𝐱_i`$ of little balls inside any single lump – by our assumption of Euclidean congruence, it doesn’t matter which one – and ranges over all values of $`i`$ and $`j`$ from $`1`$ to $`n`$ except that it is to be understood that, in this final product, in the cases $`i=j`$ the denominator is to be replaced by $`r^2`$. Clearly another way of saying exactly the same thing as this is that $`_a^{\text{grainy}}`$ is approximately equal to $`_a`$, as given by the formula $`(\text{1})`$ except that the proviso that, in the cases $`I=J`$, the terms in the denominator $`\|𝐱_I𝐱_J\|\|𝐱_I^{}𝐱_J^{}\|`$ are to be replaced by $`R^2`$ should be replaced by the proviso that, in the cases $`I=J`$, the terms in the denominator $`\|𝐱_I𝐱_J\|\|𝐱_I^{}𝐱_J^{}\|`$ are to be replaced by $`R_{\text{eff}}^2`$ where $`R_{\text{eff}}`$ is defined by $$R_{\text{eff}}=\left(\underset{i}{}\underset{j}{}\|𝐱_i𝐱_j\|\right)^{\frac{1}{n^2}}$$ where the products each go from $`1`$ to $`n`$ and it is to be understood that, in cases where $`i=j`$, the terms $`\|𝐱_i𝐱_j\|`$ are to be replaced by $`r`$. Equivalently, we may write $$R_{\text{eff}}=\left(r^n\underset{i<j}{}\|𝐱_i𝐱_j\|^2\right)^{\frac{1}{n^2}}$$ (4) (where the product is now over all $`i`$ and $`j`$ from $`1`$ to $`n`$ satisfying $`i<j`$). Thus the first part of our hope (i.e. that replacing our uniform density balls by our grainy lumps can be well-approximated by replacing $`R`$ by a suitable $`R_{\text{eff}}`$) is fulfilled. We now turn to discussing the second part of our hope, namely whether, and/or under what circumstances, $`R_{\text{eff}}`$ will turn out to be of the same “order of magnitude” as the diameter of our lump. We have not been able to answer this question in general but, instead, content ourselves with analysing one particularly simple special case of our model. Namely, where each lump is a cube consisting of $`n=(2\mathrm{}+1)^3`$ small balls of radius $`r`$ with centres at the vertices of a cubic lattice of spacing $`a`$ and side $`2\mathrm{}a`$. To spell out precisely what we mean and at the same time set up a useful notation, we assume the ball-centres of any one of these cubes to be coordinatizable so that they lie at the positions $`(i_1a,i_2a,i_3a)`$ where $`i_1`$, $`i_2`$, and $`i_3`$ are integers which each range between $`1`$ and $`2\mathrm{}+1`$. Clearly the total number of balls, $`n`$, in the cube will be $`(2\mathrm{}+1)^3`$ and we note that there will be a ball at the centre (with coordinates $`((\mathrm{}+1)a,(\mathrm{}+1)a,(\mathrm{}+1)a)`$). In line with $`(\text{2})`$, we assume $$ra\text{and}R\text{ any inter-lump spacing}.$$ (Of course, we require $`r<a/2`$ for the balls to fit into the lattice at all!) For such a lump, $`(\text{4})`$ can clearly be rewritten in the form $$R_{\text{eff}}=\left(\frac{r}{a}\right)^{1/n}\left(\underset{i<j}{}\left\|\frac{𝐱_i𝐱_j}{a}\right\|^2\right)^{\frac{1}{n^2}}a$$ (5) (here we continue to assume some system has been adopted for numbering our balls from $`1`$ to $`n=(2\mathrm{}+1)^3`$) and we notice that, for many values of the pair $`(r/a,n)`$, the prefactor, $`(r/a)^{1/n}`$ will itself be of order 1. This will be the case if, as is relevant to a ball or lump of “ordinary matter”, we take $`r/a`$ to be of the order of the ratio of the radius of the proton to the Bohr radius (i.e. around $`10^5`$) even for $`\mathrm{}=1`$ ($`n=27`$) and in fact, if $`\mathrm{}`$ is $`2`$ ($`n=125`$) or more (and it will be much more if we are taking our cube to be a model for a macroscopic lump of ordinary matter as considered in ) then this prefactor will in fact be a number close in value to 1. On the assumption that this prefactor is of order 1 (or very close to 1) our question then reduces to the question whether the quantity $$\left(\underset{i<j}{}\left\|\frac{𝐱_i𝐱_j}{a}\right\|^2\right)^{\frac{1}{n^2}}$$ is of the same order as $`\mathrm{}`$. Choosing units such that $`a=1`$, this amounts to a question about the quantity<sup>3</sup><sup>3</sup>3It is necessary to bear in mind, here and elsewhere, that $`n`$ is shorthand for $`(2\mathrm{}+1)^3`$. $$Q(\mathrm{})=\left(\underset{i<j}{}\left\|𝐱_i𝐱_j\right\|\right)^{\frac{2}{n^2}}$$ (6) where the product is over all pairs of distinct points of the cubic lattice of points with integer coordinates $`(i_1,i_2,i_3)`$ where $`i_1`$, $`i_2`$ and $`i_3`$ range between $`1`$ and $`2\mathrm{}+1`$ (from now on we shall often simply call such a finite cubic lattice of points a “cube”) and, again, we assume that some system for numbering these points from $`1`$ to $`n=(2\mathrm{}+1)^3`$ has been adopted. We formulate this question in a precise way by asking whether one can find numbers, $`c_1`$, $`c_2`$, of order 1 such that $$c_1\mathrm{}Q(\mathrm{})c_2\mathrm{}.$$ (7) It is easy to see that one can satisfy the upper bound by choosing $`c_2=2\sqrt{3}`$. To see this, notice that each term in the product in $`(\text{6})`$ is obviously less than or equal to $`2\sqrt{3}\mathrm{}`$ since this is the length of the body-diagonal of the cubic lattice and moreover there are $`n(n1)/2`$ terms in the product. We thus have that $`Q(\mathrm{})(2\sqrt{3}\mathrm{})^{n(n1)/n^2}`$ which, since $`\mathrm{}1`$, is clearly less than $`2\sqrt{3}\mathrm{}`$. However, as far as we can see, to show that there exists a (non-zero) positive $`c_1`$ such that $`c_1\mathrm{}`$ is a lower bound is not completely trivial because $`Q(\mathrm{})`$ is the $`2/n^2`$ power (i.e. $`2/(2\mathrm{}+1)^3`$ power) of a product of numbers which range in magnitude from 1 to numbers of order $`\mathrm{}`$ (the largest of the numbers will of course be $`2\sqrt{3}\mathrm{}`$). Nor, to our knowledge, does it easily follow from any existing standard mathematical result<sup>4</sup><sup>4</sup>4As far as we are aware, the notion of closest relevance to this question in the mathematical literature is the notion of the transfinite diameter of a closed bounded region (which one shows, see again , is equivalent to another notion known as the capacity of that region) of the (complex) plane. This is defined to be the supremum over all finite sets of points in the region of the $`1/n(n1)`$ power of the product of all the distances between pairs of distinct points of the set. (One can show, for example, that the transfinite diameter of a disk is equal to its radius.) If we allow ourselves to generalize this concept of transfinite diameter to closed bounded regions of three-dimensional Euclidean space, then our $`Q(\mathrm{})`$ may be seen to resemble one of the terms in the supremum which would define the transfinite diameter of a cube of side $`2\mathrm{}+1`$ (except that the $`1/n^2`$ power is taken, instead of the $`1/n(n1)`$ power). However, our interest is in $`Q(\mathrm{})`$ itself and no supremum is to be taken. However, we have succeeded in finding the (albeit probably not best-possible) number $`c_1=e^{1/3}`$ for which we can prove that it does hold: ###### Proposition. $$e^{1/3}\mathrm{}Q(\mathrm{})2\sqrt{3}\mathrm{}.$$ (Actually, one can show , by similar methods to those in the proof of the lower bound below, that if we omit $`\mathrm{}=1`$ the upper bound of $`2\sqrt{3}\mathrm{}`$ can be reduced to $`2.61\mathrm{}`$.) We present our proof of the (lower bound in) the proposition in a separate section below. Putting this proposition together with $`(\text{5})`$, we thus have the following theorem. ###### Theorem. For a lump consisting of $`n=(2\mathrm{}+1)^3`$ small balls of mass $`m`$ and radius $`r`$, centred at the vertices of a cubical region (as specified above) of side $`2\mathrm{}a`$ of a cubic lattice of spacing $`a`$, $$e^{1/3}(r/a)^{1/n}\mathrm{}aR_{\text{eff}}2\sqrt{3}(r/a)^{1/n}\mathrm{}a.$$ (Here, we recall that $`a`$ is of course assumed to be greater than $`2r`$ and $`R_{\text{eff}}`$, when inserted in place of $`R`$ in $`(\text{1})`$, is expected to give a good approximation to the multiplicative factor $`M_a^{\text{grainy}}`$ of $`(\text{3})`$ if $`a2r`$.) We remark that it is easy to calculate $`Q(\mathrm{})`$ numerically for small values of $`\mathrm{}`$ (say for $`\mathrm{}=1\mathrm{}10`$) and the numerical evidence suggests that, for such small values of $`\mathrm{}`$, $`Q(\mathrm{})`$ is well-approximated by a formula of form $`Q(\mathrm{})A\mathrm{}+B`$, i.e. $$Q(\mathrm{})/\mathrm{}A+B/\mathrm{}$$ (8) where $`A1.2`$ and $`B0.6`$. If it could be proven that there exist exact values for $`A`$ and $`B`$ such that this approximate formula holds for all $`\mathrm{}`$ with an error term which tends to zero as $`\mathrm{}`$ tends to infinity, then this would of course be a stronger result than our proposition above and would, in particular, tell us that our value for $`c_1`$, $`e^{1/3}0.71`$, can be improved to $`1.2`$. However, we have been unable to prove this. ## 2 Proof of Proposition We have already proven the upper bound above, so it remains to prove the lower bound. For this proof, we find it useful to write $`Q(\mathrm{})`$ as the product $$Q(\mathrm{})=\underset{i=1}{\overset{n}{}}\omega _i(n)$$ where $$\omega _i(n)=\underset{ji}{}\left(\left\|𝐱_i𝐱_j\right\|\right)^{\frac{1}{n^2}}$$ where the product is over all $`j`$ from $`1`$ to $`n`$ except that the factor with $`j=i`$ is omitted. We shall also assume that the numbering of the vertices from $`i=1`$ to $`n`$ is such that the centre vertex (with coordinates, as introduced above, $`(\mathrm{},\mathrm{},\mathrm{})`$) is numbered $`i=1`$. ###### Lemma. $$\omega _1\omega _j(j=1,\mathrm{},n).$$ ###### Proof. In terms of our coordinatization of our cube, we need to show that the product $$\mu _{(j_1,j_2,j_3)}=\underset{(i_1,i_2,i_3)(j_1,j_2,j_3)}{}\|(j_1,j_2,j_3)(i_1,i_2,i_3)\|$$ over all distances from an arbitrary lattice point $`(j_1,j_2,j_3)`$ to all the other lattice points in our cube is greater than or equal to the product $$\mu _{(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)}=\underset{(i_1,i_2,i_3)(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)}{}\|(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)(i_1,i_2,i_3)\|$$ over all distances from the centre point, with coordinates $`(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)`$, to all the other lattice points in our cube. (To explain the notation here, $`\omega _j=(\mu _{(j_1,j_2,j_3)})^{1/n^2}`$ if $`(j_1,j_2,j_3)`$ are the coordinates of the point numbered $`j`$. The inequality of the lemma will then follow by taking the $`1/n^2`$ power of each side of the above inequality.) We will prove this by exhibiting, for each $`(j_1,j_2,j_3)`$, a one-one onto mapping $`f_{(j_1,j_2,j_3)}`$ from (coordinatized) lattice points of our cube to themselves with the properties $$f_{(j_1,j_2,j_3)}((j_1,j_2,j_3))=(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)$$ (9) and $$\|f_{(j_1,j_2,j_3)}((i_1,i_2,i_3))(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)\|\|(i_1,i_2,i_3)(j_1,j_2,j_3)\|$$ (10) In fact, given such a mapping we will immediately have $`\mu _{(j_1,j_2,j_3)}={\displaystyle \underset{(i_1,i_2,i_3)(j_1,j_2,j_3)}{}}\|(j_1,j_2,j_3)(i_1,i_2,i_3)\|`$ $`{\displaystyle \underset{(i_1,i_2,i_3)(j_1,j_2,j_3)}{}}\|f_{(j_1,j_2,j_3)}((j_1,j_2,j_3))f_{(j_1,j_2,j_3)}((i_1,i_2,i_3))\|`$ $`={\displaystyle \underset{(i_1,i_2,i_3)(j_1,j_2,j_3)}{}}\|(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)f_{(j_1,j_2,j_3)}((i_1,i_2,i_3))\|`$ $`={\displaystyle \underset{(\widehat{i}_1,\widehat{i}_2,\widehat{i}_3)(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)}{}}\|(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)(\widehat{i}_1,\widehat{i}_2,\widehat{i}_3)\|`$ $`=\mu _{(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)}`$ where, in the penultimate equality, we have defined $`(\widehat{i}_1,\widehat{i}_2,\widehat{i}_3)`$ to be $`f_{(j_1,j_2,j_3)}((i_1,i_2,i_3))`$. It therefore remains to exhibit a mapping, $`f_{(j_1,j_2,j_3)}`$, with the above properties. To do this, we define $$f_{(j_1,j_2,j_3)}((i_1,i_2,i_3))=(\mathrm{}+1j_1+i_1,\mathrm{}+1j_2+i_2,\mathrm{}+1j_3+i_3)(\text{mod}2\mathrm{}+1)$$ where we use a non-conventional notion of addition modulo $`q`$ in which $`q(\text{mod}q)`$ is deemed to be $`q`$ rather than $`0`$. One may understand this definition geometrically as a rigid translation of all the points of our cube where, however, if the would-be destination of a point is outside of our cube, the point instead gets mapped to the counterpart-position in our cubic lattice to the position it would arrive at in a neighbouring cube, were our cube to be surrounded by similar neighbours in a cubic lattice arrangement. Clearly this is a bijection and satisfies (9). To show it satisfies (10) we calculate, using the obvious properties of our notion of “mod”: $`\|f_{(j_1,j_2,j_3)}((i_1,i_2,i_3))(\mathrm{}+1,\mathrm{}+1,\mathrm{}+1)\|=([((\mathrm{}+1j_1+i_1)(\text{mod}2\mathrm{}+1))\mathrm{}1]^2`$ $`+\left[((\mathrm{}+1j_2+i_2)(\text{mod}2\mathrm{}+1))\mathrm{}1\right]^2`$ $`+[((\mathrm{}+1j_3+i_3)(\text{mod}2\mathrm{}+1))\mathrm{}1]^2)^{1/2}`$ $`\left((j_1i_1)^2+(j_2i_2)^2+(j_3i_3)^2\right)^{1/2}=\|(j_1,j_2,j_3)(i_1,i_2,i_3)\|.`$ Before proceeding with the proofs of the two bounds in our proposition, it will be useful to define, for a given cube (i.e. cubic lattice of points) of side $`2\mathrm{}+1`$, certain special sets of lattice points. First, we define the “$`k`$th centre shell” ($`k\mathrm{}`$) to be the set of lattice points on the surface of the cube of side $`2k+1`$ centred on the centre of our given cube (so that a cube of side $`2\mathrm{}+1`$ would have a total of $`\mathrm{}`$ centre shells). ###### Proposition (Lower Bound). For a cube of side $`2\mathrm{}+1`$, we have $$\frac{Q(\mathrm{})}{\mathrm{}}e^{\frac{1}{3}}\mathrm{}$$ (12) ###### Proof. For such a cube we have that $$n=(2\mathrm{}+1)^3$$ Let us denote the number of lattice points in the $`k`$th centre shell by $`n_k`$. This is given by $$n_k=6(2k+1)^212(2k+1)+8$$ We first use the lemma to estimate $$Q(\mathrm{})=\underset{i=1}{\overset{n}{}}\omega _i(n)\omega _1(n)^n$$ We can rewrite $`\omega _1(n)`$ by a relabelling of the lattice points $$\omega _1(n)=\left(\underset{i=1}{\overset{n}{}}\|𝐱_1𝐱_i\|\right)^{\frac{1}{n^2}}=\left(\underset{k=1}{\overset{\mathrm{}}{}}\underset{d=1}{\overset{n_k}{}}\|𝐱_1𝐱_{d_k}\|\right)^{\frac{1}{n^2}}$$ (13) where $`𝐱_{d_k}`$ is the coordinate-triple of the $`d^{th}`$ lattice point in the $`k`$th centre shell. Each lattice point in the $`k`$th centre shell is at least a distance $`k`$ away from the centre lattice point so we have $$Q(\mathrm{})\alpha (\mathrm{})^{\frac{1}{n}}$$ (14) where $$\alpha (\mathrm{})=\underset{k=1}{\overset{\mathrm{}}{}}k^{n_k}$$ (15) From (14), it can be seen that to prove (12) it will be sufficient to prove $$\alpha (\mathrm{})\left(e^{\frac{1}{3}}\mathrm{}\right)^n=\left(e^{\frac{1}{3}}\mathrm{}\right)^{(2\mathrm{}+1)^3}$$ (16) This can be proved by induction on $`\mathrm{}`$. For all $`\mathrm{}`$, let $`P(\mathrm{})`$ be the proposition that (16) is true. We have $$\alpha (1)=1e^{\frac{1}{3}},$$ so $`P(1)`$ is true. So, if we can show $`P(\mathrm{})P(\mathrm{}+1)`$ ($`\mathrm{}`$) then the proposition will have been proved to be true. From (15) and (16) $$\alpha (\mathrm{}+1)=\alpha (\mathrm{})(\mathrm{}+1)^{n_{\mathrm{}+1}}\left(e^{\frac{1}{3}}\mathrm{}\right)^n(\mathrm{}+1)^{n_{\mathrm{}+1}}$$ Thus to now prove our proposition it will be sufficient to prove that $$\left(e^{\frac{1}{3}}\mathrm{}\right)^n(\mathrm{}+1)^{n_{\mathrm{}+1}}\left(e^{\frac{1}{3}}(\mathrm{}+1)\right)^{(2\mathrm{}+3)^3}$$ This is easily seen to be equivalent to $$e^{\frac{1}{3}}\left(1+\frac{1}{\mathrm{}}\right)^{\frac{1}{3}\mathrm{}+\left(\frac{132\mathrm{}^2+72\mathrm{}1}{24\mathrm{}^2}\right)}.$$ However the above equation clearly holds so we have proved our inductive proposition (and hence the lower bound in our main proposition). We thank an anonymous referee for pointing out to us that numerical evidence suggests the linear form for $`Q(\mathrm{})`$ of equation $`(\text{8})`$. VA is grateful to Richard Hunter for useful discussions about, and constructive criticism of, some of the mathematical proofs in this paper. VA is also grateful to Jim Brink for advice on numerical computation. VA also thanks PPARC for a research studentship. BSK is grateful to Richard Hall for telling him about the concept of transfinite diameter ($`=`$ capacity) of regions of the plane. BSK would also like to thank the Leverhulme foundation for a Leverhulme fellowship (RF&G/9/RFG/2002/0377) from October 2002 to June 2003 during the course of which this research was begun. ## References
warning/0506/nlin0506043.html	ar5iv	text	# Chaotic Dynamics of N–degree of Freedom Hamiltonian Systems ## 1 Introduction Chaotic behavior in Hamiltonian systems with many degrees of freedom has been the subject of intense investigation in the last fifty years, see e.g. \[Lichtenberg & Lieberman, 1991; MacKay & Meiss, 1987; Wiggins, 1988\] and Simó ed., . By degrees of freedom (dof) we are referring to the number of canonically conjugate pairs of positions and momentum variables, $`q_k`$ and $`p_k`$ respectively, with $`k=1,2,\mathrm{},N`$. The relevance of these systems to problems of practical concern cannot be overemphasized. Their applications range from the stability of the solar system Contopoulos, and the containment of charged particles in high intensity magnetic fields Lichtenberg & Lieberman, to the blow–up of hadron beams in high energy accelerators Scandale & Turchetti, eds. and the understanding of the properties of simple molecules and hydrogen–bonded systems Bountis ed., ; Prosmiti & Farantos, . One of the most fundamental areas in which the dynamics of multi–degree of freedom Hamiltonian systems has played (and continues to play) a crucial role is the study of transport phenomena in one–dimensional ($`1`$D) lattices and the role of chaos in providing a link between deterministic and statistical behavior Chirikov, ; Lichtenberg & Lieberman, ; Ford, . In this context, a lattice of $`N`$ dof is expected, in the thermodynamic limit ($`N\mathrm{}`$ at fixed energy density $`E/N`$) to exhibit chaotic behavior for almost all initial configurations, satisfying at least the property of ergodicity. This would allow the use of probability densities, leading from the computation of orbits to the study of statistical quantities like ensemble averages and transport coefficients. Chaotic regions, where nearby solutions diverge exponentially from each other, provide an excellent “stage” on which such a desired transition from classical to statistical mechanics can occur. However, the presence of significantly sized islands (or tori) of quasiperiodic motion, in which the dynamics is “stable” for long times, preclude the success of this scenario and make any attempt at a globally valid statistical description seriously questionable. These tori occur e.g. around simple periodic orbits which are stable under small perturbations and, if their size does not shrink to zero as $`N`$ or $`E`$ increases, their presence can attribute global consequences to a truly local phenomenon. That such important periodic orbits do exist in Hamiltonian lattices, even in the $`N\mathrm{}`$ limit, could not have been more dramatically manifested than in the remarkable discovery of discrete breathers (see e.g. Flach & Willis, ), which by now have been observed in a great many experimental situations \[Eisenberg et al., 1998; Schwarz et al., 1999; Fleischer et al., 2003; Sato et al., 2003\]. Discrete breathers are precisely one such kind of stable periodic orbits, which also happen to be localized in space, thus representing a very serious limitation to energy transport in nonlinear lattices. Then, there were, of course, the famous numerical experiments of Fermi, Pasta and Ulam of the middle $`1950`$’s Fermi et al., , which demonstrated the existence of recurrences that prevent energy equipartition among the modes of certain $`1`$D lattices, containing nonlinear interactions between nearest neighbors. These finite, so–called FPU Hamiltonian systems were later shown to exhibit a transition to “global” chaos, at high enough energies where major resonances overlap Izrailev & Chirikov, . Before that transition, however, an energy threshold to a “weak” form of chaos was later discovered that relies on the interaction of the first few lowest frequency modes and, at least for the FPU system, does appear to ensure equipartition among all modes \[De Luca et al., $`1995`$; De Luca & Lichtenberg, $`2002`$\]. Interestingly enough, very recently, this transition to “weak” chaos was shown to be closely related to the destabilization of one of the lowest frequency nonlinear normal mode of this FPU system Flach et al., . Thus, today, $`50`$ years after its famous discovery, the Fermi–Pasta–Ulam problem and its transition from recurrences to true statistical behavior is still a subject of ongoing investigation Berman & Izrailev, . In this paper, we have sought to approach the problem of global chaos in Hamiltonian systems, by considering two paradigms of $`N`$ dof, $`1`$D nonlinear lattices, with very different origins. One is the famous FPU lattice mentioned above, with quadratic and quartic nearest neighbor interactions, described by the Hamiltonian $$H=\frac{1}{2}\underset{j=1}{\overset{N}{}}\dot{x}_j^2+\underset{j=0}{\overset{N}{}}\left(\frac{1}{2}(x_{j+1}x_j)^2+\frac{1}{4}\beta (x_{j+1}x_j)^4\right)=E$$ (1) where $`x_j`$ is the displacement of the $`j`$th particle from its equilibrium position, $`\dot{x}_j`$ is the corresponding canonically conjugate momentum of $`x_j`$, $`\beta `$ is a positive real constant and $`E`$ is the value of the Hamiltonian representing the total energy of the system. The other one is obtained by a discretization of a partial differential equation (PDE) of the nonlinear Schrödinger type referred to as the Gross–Pitaevskii equation Dalfovo et al., , which in dimensionless form reads $$i\mathrm{}\frac{\mathrm{\Psi }(x,t)}{t}=\frac{\mathrm{}^2}{2}\frac{^2\mathrm{\Psi }(x,t)}{x^2}+V(x)\mathrm{\Psi }(x,t)+g\|\mathrm{\Psi }(x,t)\|^2\mathrm{\Psi }(x,t),i^2=1$$ (2) where $`\mathrm{}`$ is the Planck constant, $`g`$ is a positive constant (repulsive interactions between atoms in the condensate) and $`V(x)`$ is an external potential. Equation (2) is related to the phenomenon of Bose–Einstein Condensation (BEC) Ketterle et al., . Here we consider the simple case $`V(x)=0`$, $`\mathrm{}=1`$ and discretize the $`x`$–dependence of the complex variable $`\mathrm{\Psi }(x,t)\mathrm{\Psi }_j(t)`$ in (2), approximating the second order derivative by $`\mathrm{\Psi }_{xx}\frac{\mathrm{\Psi }_{j+1}+\mathrm{\Psi }_{j1}2\mathrm{\Psi }_j}{\delta x^2}`$. Setting then $`\mathrm{\Psi }_j(t)=q_j(t)+ip_j(t),i^2=1,j=1,2,\mathrm{},N`$ and $`\|\mathrm{\Psi }(x,t)\|^2=q_j^2(t)+p_j^2(t)`$, one immediately obtains from the above PDE (2) a set of ordinary differential equations (ODEs) for the canonically conjugate variables, $`p_j`$ and $`q_j`$, described by the BEC Hamiltonian Trombettoni & Smerzi, ; Smerzi & Trombettoni, $$H=\frac{1}{2}\underset{j=1}{\overset{N}{}}(p_j^2+q_j^2)+\frac{\gamma }{8}\underset{j=1}{\overset{N}{}}(p_j^2+q_j^2)^2\frac{ϵ}{2}\underset{j=1}{\overset{N}{}}(p_jp_{j+1}+q_jq_{j+1})=E$$ (3) where $`\gamma >0`$ and $`ϵ=1`$ are constant parameters, $`g=\frac{\gamma }{2}>0`$ with $`\delta x=1`$ and $`E`$ is the total energy of the system. In Sec. 2, we study these Hamiltonians, focusing on some simple periodic orbits (SPOs), which are known in closed form and whose local (linear) stability analysis can be carried out to arbitrary accuracy. By SPOs, we refer here to periodic solutions where all variables return to their initial state after only one maximum and one minimum in their oscillation, i.e. all characteristic frequencies have unit ratios. In particular, we examine first their bifurcation properties to determine whether they remain stable for arbitrarily large $`E`$ and $`N`$, having perhaps finitely sized islands of regular motion around them. This was found to be true only for the so–called in–phase–mode (IPM) of the BEC Hamiltonian (3). SPOs corresponding to out–of–phase motion (OPM) of either the FPU (1) or the BEC system (3) destabilize at energy densities $`\frac{E_c}{N}N^\alpha `$, with $`\alpha =1`$ or $`2`$ (for the SPOs we studied), as $`N\mathrm{}`$. The same result was also obtained for what we call the OHS mode of the FPU system Ooyama et al., , where all even indexed particles are stationary and all others execute out–of–phase oscillation, under fixed or periodic boundary conditions. All these are in agreement with detailed analytical results obtained for families of SPOs of the same FPU system under periodic boundary conditions (see e.g. Poggi & Ruffo, ). Then, in Sec. 3, we vary the values of $`E`$ and $`N`$ and study the behavior of the Lyapunov exponents of the OPMs of the FPU and BEC Hamiltonians. We find that, as the eigenvalues of the monodromy matrix exit the unit circle on the real axis, at energies $`0<E_cE_1<E_2<\mathrm{}`$, for fixed $`N`$, all positive Lyapunov exponents $`L_i,i=1,\mathrm{},N1`$, increase following two distinct power laws, $`L_iE^{B_k}`$, $`B_k>0,k=1,2`$, with the $`B_k`$’s as reported in the literature, Rechester et al., ; Benettin, and Livi et al., . Furthermore, as the energy $`E`$ grows at fixed $`N`$, the real eigenvalues of the OPM orbit of FPU continue to move away from $`1`$, unlike the OPM of the BEC Hamiltonian, where for very large $`E`$ all these eigenvalues tend to return to $`+1`$. Interestingly enough, the IPM of the BEC Hamiltonian remains stable for all the energies and number of dof we studied! In Sec. 4, we turn to the question of the “size” of islands around stable SPOs and use the Smaller Alignment Index (SALI), introduced in earlier papers Skokos, ; Skokos et al., 2003a ; Skokos et al., 2003b ; Skokos et al., to distinguish between regular and chaotic trajectories in our two Hamiltonians. First, we verify again in these multi–degree of freedom systems the validity of the SALI dependence on the two largest Lyapunov exponents $`L_1`$ and $`L_2`$ in the case of chaotic motion, to which it owes its effectiveness and predictive power. Then, computing the SALI, at points further and further away from stable SPOs, we determine approximately the “magnitude” of these islands and find that it vanishes (as expected) at the points where the corresponding OPMs destabilize. In fact, for the OPM of the FPU system the size of the islands decreases monotonically before destabilization while for the BEC orbit the opposite happens! Even more remarkably, for the IPM of the BEC Hamiltonian (3), not only does the “size” of the island not vanish, it even grows with increasing energy and remains of considerable magnitude for all the values of $`E`$ and $`N`$ we considered. Finally, in Sec. 5, using as initial conditions the unstable SPOs, we compute the Lyapunov spectra of the FPU and BEC systems in the so–called thermodynamic limit, i.e. as the energy $`E`$ and the number of dof $`N`$ grow indefinitely, with energy density $`\frac{E}{N}`$ fixed. First, we find that Lyapunov exponents, fall on smooth curves of the form $`L_iL_1e^{\alpha i/N}`$, for both systems. Then, computing the Kolmogorov–Sinai entropy $`h_{KS}`$, as the sum of the positive Lyapunov exponents Pesin, ; Hilborn, , in the thermodynamic limit, we find, for both Hamiltonians, that $`h_{KS}`$ is an extensive quantity as it grows linearly with $`N`$, demonstrating that in their chaotic regions the FPU and BEC Hamiltonians behave as ergodic systems of statistical mechanics. ## 2 Simple Periodic Orbits (SPOs) and Local Stability Analysis ### 2.1 The FPU model #### 2.1.1 Analytical expressions of the OHS mode We consider a $`1`$D lattice of $`N`$ particles with nearest neighbor interactions given by the Hamiltonian Fermi et al., $$H=\frac{1}{2}\underset{j=1}{\overset{N}{}}\dot{x}_j^2+\underset{j=0}{\overset{N}{}}\left(\frac{1}{2}(x_{j+1}x_j)^2+\frac{1}{4}\beta (x_{j+1}x_j)^4\right)$$ (4) where $`x_j`$ is the displacement of the $`j`$th particle from its equilibrium position, $`\dot{x}_j`$ is the corresponding canonically conjugate momentum of $`x_j`$ and $`\beta `$ is a positive real constant. Imposing fixed boundary conditions $$x_0(t)=x_{N+1}(t)=0,t$$ (5) one finds a simple periodic orbit first studied by Ooyama et al., , taking $`N`$ odd, which we shall call the OHS mode (using Ooyama’s, Hirooka’s and Saitô’s initials) $$\widehat{x}_{2j}(t)=0,\widehat{x}_{2j1}(t)=\widehat{x}_{2j+1}(t)\widehat{x}(t),j=1,\mathrm{},\frac{N1}{2}.$$ (6) Here, we shall examine analytically the stability properties of this mode and determine the energy range $`0EE_c(N)`$ over which it is linearly stable. The equations of motion associated with Hamiltonian (4) are $$\ddot{x}_j(t)=x_{j+1}2x_j+x_{j1}+\beta \left((x_{j+1}x_j)^3(x_jx_{j1})^3\right),j=1,\mathrm{},N$$ (7) whence, using the boundary condition (5) and the expressions (6) for every $`j=1,3,\mathrm{},N2,N`$, we arrive at a single equation $$\ddot{\widehat{x}}(t)=2\widehat{x}(t)2\beta \widehat{x}^3(t)$$ (8) describing the anharmonic oscillations of all odd particles of the initial lattice. The solution of (8) is, of course, well–known in terms of Jacobi elliptic functions Abramowitz & Stegun, and can be written as $$\widehat{x}(t)=𝒞cn(\lambda t,\kappa ^2)$$ (9) where $$𝒞^2=\frac{2\kappa ^2}{\beta (12\kappa ^2)},\lambda ^2=\frac{2}{12\kappa ^2}$$ (10) and $`\kappa ^2`$ is the modulus of the $`cn`$ elliptic function. The energy per particle of this SPO is then given by $$\frac{E}{N+1}=\frac{1}{4}𝒞^2(2+𝒞^2\beta )=\frac{\kappa ^2\kappa ^4}{(12\kappa ^2)^2\beta }.$$ (11) #### 2.1.2 Stability analysis of the OHS mode Setting $`x_j=\widehat{x}_j+y_j`$ in (7) and keeping up to linear terms in $`y_j`$ we get the corresponding variational equations for the OHS mode (6) $$\ddot{y}_j=(1+3\beta \widehat{x}^2)(y_{j1}2y_j+y_{j+1}),j=1,\mathrm{},N$$ (12) where $`y_0=y_{N+1}=0`$. Using the standard method of diagonalization of linear algebra, we can separate these variational equations to $`N`$ uncoupled independent Lamé equations \[Abramowitz & Stegun, 1965\] $$\ddot{z}_j(t)+4(1+3\beta \widehat{x}^2)\mathrm{sin}^2\left(\frac{\pi j}{2(N+1)}\right)z_j(t)=0,j=1,\mathrm{},N$$ (13) where the $`z_j`$ variations are simple linear combinations of $`y_j`$’s. Using (9) and changing variables to $`u=\lambda t`$, Eq. (13) takes the form $$z_j^{\prime \prime }(u)+2\left(1+4\kappa ^26\kappa ^2\stackrel{2}{sn}(u,\kappa ^2)\right)\mathrm{sin}^2\left(\frac{\pi j}{2(N+1)}\right)z_j(u)=0,j=1,\mathrm{},N$$ (14) where we have used the identity $`cn^2(u,\kappa ^2)=1sn^2(u,\kappa ^2)`$ and primes denote differentiation with respect to $`u`$. Equation (14) is an example of Hill’s equation \[Copson, $`1935`$; Magnus & Winkler, $`1966`$\] $$z^{\prime \prime }(u)+Q(u)z(u)=0$$ (15) where $`Q(u)`$ is a $`T`$–periodic function ($`Q(u)=Q(u+T)`$) with $`T=2𝒦`$ and $`𝒦𝒦(\kappa ^2)`$ is the elliptic integral of the first kind. According to Floquet theory Magnus & Winkler, the solutions of Eq. (15) are bounded (or unbounded) depending on whether the Floquet exponent $`\alpha `$, given by $$\mathrm{cos}\left(2\alpha 𝒦(\kappa ^2)\right)=12\mathrm{sin}^2\left(𝒦(\kappa ^2)\sqrt{a_0}\right)det\left(𝐃(0)\right)$$ (16) is real (or imaginary). The matrix $`𝐃(\alpha )`$ is called Hill’s matrix and in our case its entries are given by $$[𝐃(\alpha )]_{n,m}\frac{a_{nm}}{a_0\left(\alpha +\frac{n\pi }{𝒦(\kappa ^2)}\right)^2}+\delta _{n,m}$$ (17) where $`\delta _{n,m}=\{\begin{array}{cc}1,\hfill & n=m\hfill \\ 0,\hfill & nm\hfill \end{array}`$, is the Kronecker delta with $`n`$, $`m`$ and the $`a_n`$’s are the coefficients of the Fourier series expansion of $`Q(u)`$, $$Q(u)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a_ne^{\frac{in\pi u}{𝒦(\kappa ^2)}}.$$ (18) Thus, Eq. (16) gives a stability criterion for the OHS mode (6), by the condition $$\left\|12\mathrm{sin}^2\left(𝒦(\kappa ^2)\sqrt{a_0}\right)det\left(𝐃(0)\right)\right\|=\{\begin{array}{cc}<1,\hfill & \mathrm{stable}\mathrm{mode}\hfill \\ >1,\hfill & \mathrm{unstable}\mathrm{mode}\hfill \end{array}.$$ (19) In this case, the Fourier coefficients of Hill’s matrix $`𝐃(0)`$ are given by the relations \[Copson, $`1935`$\] $`a_0`$ $`=`$ $`2(5+4\kappa ^2+6{\displaystyle \frac{(\kappa ^2)}{𝒦(\kappa ^2)}})\mathrm{sin}^2\left({\displaystyle \frac{\pi j}{2(N+1)}}\right)`$ (20) $`a_n`$ $`=`$ $`2{\displaystyle \frac{6n\pi ^2q^n}{𝒦^2(\kappa ^2)(1q^{2n})}}\mathrm{sin}^2({\displaystyle \frac{\pi j}{2(N+1)}}),n0`$ (21) where $`qe^{\pi \frac{𝒦^{}}{𝒦}}`$, $`𝒦𝒦(\kappa ^2)`$ and $`(\kappa ^2)`$ are the elliptic integrals of the first and second kind respectively and $`𝒦^{}𝒦(\kappa ^2)`$ with $`\kappa ^21\kappa ^2`$. In the evaluation of $`𝐃(0)`$ in Eq. (19) we have used $`121`$ terms in the Fourier series expansion of $`sn^2(u,\kappa ^2)`$ (that is, $`121\times 121`$ Hill’s determinants of $`𝐃(0)`$). Thus, we determine with accuracy $`10^8`$ the $`\kappa ^2\kappa _j^2`$ values at which the Floquet exponent $`\alpha `$ in (16) becomes zero and the $`z_j(u)`$ in (14) become unbounded. We thus find that the first variation $`z_j(u)`$ to become unbounded as $`\kappa ^2`$ increases is $`j=\frac{N1}{2}`$ and the energy values $`E_c`$ at which this happens (see Eq. (11)) are listed in Table 1 for $`\beta =1.04`$. The $`z_j(u)`$ variation with $`j=\frac{N+1}{2}`$ has Floquet exponent $`\alpha `$ equal to zero for every $`\kappa ^2[0,\frac{1}{2}]`$, that is $`z_{\frac{N+1}{2}}(u)`$ corresponds to variations along the orbit. Next we vary $`N`$ and calculate the destabilization energy per particle $`\frac{E_c}{N}`$ for $`\beta =1.04`$ at which the OHS nonlinear mode (6) becomes unstable. Plotting the results in Fig. 1, we see that $`\frac{E_c}{N}`$ decreases for large $`N`$ with a simple power–law, as $`1/N`$. #### 2.1.3 Analytical study of an OPM solution We now turn to the properties of another SPO of the FPU Hamiltonian (4), studied in Budinsky & Bountis, ; Poggi & Ruffo, and \[Cafarella et al., 2003\]. In particular, imposing the periodic boundary conditions $$x_{N+1}(t)=x_1(t),t$$ (22) with $`N`$ even, we analyze the stability properties of the out–of–phase mode (OPM) defined by $$\widehat{x}_j(t)=\widehat{x}_{j+1}(t)\widehat{x}(t),j=1,\mathrm{},N.$$ (23) In this case the equations of motion (7) reduce also to a single differential equation $$\ddot{\widehat{x}}(t)=4\widehat{x}(t)16\beta \widehat{x}^3(t)$$ (24) describing the anharmonic oscillations of all particles of the initial lattice. The solution of Eq. (24) can again be written as an elliptic $`cn`$–function $$\widehat{x}(t)=𝒞cn(\lambda t,\kappa ^2)$$ (25) with $$𝒞^2=\frac{\kappa ^2}{2\beta (12\kappa ^2)},\lambda ^2=\frac{4}{12\kappa ^2}.$$ (26) The energy per particle of the nonlinear OPM (23) is given by $$\frac{E}{N}=2𝒞^2(1+2𝒞^2\beta )=\frac{\kappa ^2\kappa ^4}{(12\kappa ^2)^2\beta }$$ (27) in this case. We study the linear stability of the OPM (23) following a similar analysis to the one performed in the case of the OHS mode of Sec. 2.1.2. In this case, the corresponding variational equations have the form $$\ddot{y}_j=(1+12\beta \widehat{x}^2)(y_{j1}2y_j+y_{j+1}),j=1,\mathrm{},N$$ (28) where $`y_1=y_{N+1}`$. After the appropriate diagonalization, the above equations are transformed to a set of $`N`$ uncoupled independent Lamé equations, which take the form $$z_j^{\prime \prime }(u)+\left(1+4\kappa ^26\kappa ^2\stackrel{2}{sn}(u,\kappa ^2)\right)\mathrm{sin}^2\left(\frac{\pi j}{N}\right)z_j(u)=0,j=1,\mathrm{},N$$ (29) after changing to the new time variable $`u=\lambda t`$. Primes denotes again differentiation with respect to $`u`$. As in Sec. 2.1.2, Eq. (19) gives a stability criterion for the nonlinear mode (23) with $`a_0`$ $`=`$ $`(5+4\kappa ^2+6{\displaystyle \frac{(\kappa ^2)}{𝒦(\kappa ^2)}})\mathrm{sin}^2({\displaystyle \frac{\pi j}{N}}),`$ (30) $`a_n`$ $`=`$ $`{\displaystyle \frac{6n\pi ^2q^n}{𝒦^2(\kappa ^2)(1q^{2n})}}\mathrm{sin}^2({\displaystyle \frac{\pi j}{N}}),n0.`$ (31) Proceeding in the same way as with the OHS mode, we find that the first variation $`z_j(u)`$ in Eq. (29) that becomes unbounded (for $`\beta =1`$ and $`N`$ even) is $`j=\frac{N}{2}1`$, in accordance with Budinsky & Bountis, , Poggi & Ruffo, , Cafarella et al., and that this occurs at the energy values $`E_c`$, listed in Table 2 for $`\beta =1`$. The $`z_j(u)`$ variation with $`j=\frac{N}{2}`$ corresponds to variations along the orbit. Taking now many values of N (even) and computing the energy per particle $`\frac{E_c}{N}`$ for $`\beta =1`$ at which the OPM (23) first becomes unstable, we plot the results in Fig. 2 and find that it also decreases following a power–law of the form $`1/N^2`$. ### 2.2 The BEC model #### 2.2.1 Analytical expressions for SPOs The Hamiltonian of the Bose–Einstein Condensate (BEC) model studied in this paper is given by $$H=\frac{1}{2}\underset{j=1}{\overset{N}{}}(p_j^2+q_j^2)+\frac{\gamma }{8}\underset{j=1}{\overset{N}{}}(p_j^2+q_j^2)^2\frac{ϵ}{2}\underset{j=1}{\overset{N}{}}(p_jp_{j+1}+q_jq_{j+1})$$ (32) where $`\gamma `$ and $`ϵ`$ are real constants, which we take here to be $`\gamma =ϵ=1`$. The Hamiltonian (32) possesses a second integral of the motion given by $$F=\underset{j=1}{\overset{N}{}}(p_j^2+q_j^2)$$ (33) and therefore chaotic behavior can only occur for $`N3`$. Imposing periodic boundary conditions to the BEC Hamiltonian (32) $`q_{N+1}(t)`$ $`=`$ $`q_1(t)\mathrm{and}`$ $`p_{N+1}(t)`$ $`=`$ $`p_1(t),t`$ (34) we analyze the stability properties of the in–phase–mode (IPM) $`q_j(t)`$ $``$ $`\widehat{q}(t),`$ $`p_j(t)`$ $``$ $`\widehat{p}(t)j=1,\mathrm{},N,N\mathrm{and}N2`$ (35) and of the out–of–phase mode (OPM) $`q_j(t)`$ $`=`$ $`q_{j+1}(t)\widehat{q}(t),`$ $`p_j(t)`$ $`=`$ $`p_{j+1}(t)\widehat{p}(t)j=1,\mathrm{},N,\mathrm{with}N\mathrm{only}\mathrm{even}`$ (36) and determine the energy range $`0EE_c(N)`$ over which these two SPOs are linearly stable. In both cases, the corresponding equations of motion $`\dot{q}_j`$ $`=`$ $`p_j+{\displaystyle \frac{\gamma }{2}}(p_j^2+q_j^2)p_j{\displaystyle \frac{ϵ}{2}}(p_{j1}+p_{j+1}),`$ $`\dot{p}_j`$ $`=`$ $`\left(q_j+{\displaystyle \frac{\gamma }{2}}(p_j^2+q_j^2)q_j{\displaystyle \frac{ϵ}{2}}(q_{j1}+q_{j+1})\right),j=1,\mathrm{},N`$ (37) give for the IPM solution $$\dot{\widehat{q}}=\widehat{p}+\frac{\gamma }{2}(\widehat{p}^2+\widehat{q}^2)\widehat{p}ϵ\widehat{p},\dot{\widehat{p}}=\left(\widehat{q}+\frac{\gamma }{2}(\widehat{p}^2+\widehat{q}^2)\widehat{q}ϵ\widehat{q}\right)$$ (38) and for the OPM $$\dot{\widehat{q}}=\widehat{p}+\frac{\gamma }{2}(\widehat{p}^2+\widehat{q}^2)\widehat{p}+ϵ\widehat{p},\dot{\widehat{p}}=\left(\widehat{q}+\frac{\gamma }{2}(\widehat{p}^2+\widehat{q}^2)\widehat{q}+ϵ\widehat{q}\right).$$ (39) From Eq. (33) we note that the second integral becomes for both SPOs $$F=N(\widehat{q}^2+\widehat{p}^2)$$ (40) yielding for the IPM solution $$\dot{\widehat{q}}=\left(1ϵ+\frac{\gamma F}{2N}\right)^2\widehat{p},\dot{\widehat{p}}=\left(1ϵ+\frac{\gamma F}{2N}\right)^2\widehat{q}$$ (41) and for the OPM $$\dot{\widehat{q}}=\left(1+ϵ+\frac{\gamma F}{2N}\right)^2\widehat{p},\dot{\widehat{p}}=\left(1+ϵ+\frac{\gamma F}{2N}\right)^2\widehat{q}.$$ (42) The above equations imply for both SPOs that their solutions are simple trigonometric functions $`\ddot{\widehat{q}}(t)`$ $`=`$ $`\omega ^2\widehat{q}(t)`$ $`\widehat{q}(t)`$ $`=`$ $`C_1\mathrm{cos}(\omega t)+C_2\mathrm{sin}(\omega t)`$ (43) $`\widehat{p}(t)`$ $`=`$ $`C_1\mathrm{sin}(\omega t)+C_2\mathrm{cos}(\omega t)`$ with $`\omega =1ϵ+\frac{\gamma F}{2N}`$ for the IPM and $`\omega =1+ϵ+\frac{\gamma F}{2N}`$ for the OPM with $`C_1`$ and $`C_2`$ real, arbitrary constants, where $`F=NA`$ and $`A=C_1^2+C_2^2`$. The energy per particle for these two orbits is then given by $$\frac{E}{N}=\frac{1ϵ}{2}A+\frac{\gamma }{8}A^2(\mathrm{IPM})\mathrm{and}\frac{E}{N}=\frac{1+ϵ}{2}A+\frac{\gamma }{8}A^2(\mathrm{OPM}).$$ (44) Such SPOs have also been studied in the case of the integrable so–called dimer problem by other authors Aubry et al., , who were interested in comparing the classical with the quantum properties of the BEC Hamiltonian (32). #### 2.2.2 Stability analysis of the SPOs Setting now $`q_j`$ $`=`$ $`\widehat{q}_j+x_j,`$ $`p_j`$ $`=`$ $`\widehat{p}_j+y_j,j=1,\mathrm{},N`$ (45) in the equations of motion (2.2.1) and keeping up to linear terms in $`x_j`$ and $`y_j`$ we get the corresponding variational equations for both SPOs $`\dot{x}_j`$ $`=`$ $`\left({\displaystyle \frac{ϵ}{2}}\right)y_{j1}+L_+y_j\left({\displaystyle \frac{ϵ}{2}}\right)y_{j+1}+Kx_j,`$ $`\dot{y}_j`$ $`=`$ $`\left(\left({\displaystyle \frac{ϵ}{2}}\right)x_{j1}+L_{}x_j\left({\displaystyle \frac{ϵ}{2}}\right)x_{j+1}+Ky_j\right),j=1,\mathrm{},N`$ (46) where $`x_0=x_N,y_0=y_N`$ and $`x_{N+1}=x_1,y_{N+1}=y_1`$ (periodic boundary conditions) and $`K`$ $`=`$ $`\gamma \widehat{p}_j\widehat{q}_j={\displaystyle \frac{1}{2}}\gamma (2C\mathrm{cos}(2\omega t)+B\mathrm{sin}(2\omega t)),`$ $`L_+`$ $`=`$ $`1+{\displaystyle \frac{\gamma }{2}}(\widehat{p}_j^2+\widehat{q}_j^2)+\gamma \widehat{p}_j^2=1+A\gamma +{\displaystyle \frac{1}{2}}B\gamma \mathrm{cos}(2\omega t)C\gamma \mathrm{sin}(2\omega t),`$ (47) $`L_{}`$ $`=`$ $`1+{\displaystyle \frac{\gamma }{2}}(\widehat{p}_j^2+\widehat{q}_j^2)+\gamma \widehat{q}_j^2=1+A\gamma {\displaystyle \frac{1}{2}}B\gamma \mathrm{cos}(2\omega t)+C\gamma \mathrm{sin}(2\omega t)`$ where $`B=C_2^2C_1^2`$ and $`C=C_1C_2`$ are real constants. Unfortunately, it is not as easy to uncouple this linear system of differential equations (2.2.2), as it was in the FPU case, in order to study analytically the linear stability of these two SPOs. We can, however, compute with arbitrarily accuracy for every given $`N`$ the eigenvalues of the corresponding monodromy matrix of the IPM and OPM solutions of the BEC Hamiltonian (32). Thus, in the case of the OPM (2.2.1) we computed for some even values of $`N`$ the energy thresholds $`E_c`$ at which this SPO becomes unstable (see Table 3). Plotting the results in Fig. 3 we observe again that $`E_c/N`$ decreases with $`N`$ following a power–law $`1/N^2`$ as in the case of the OPM of the FPU model. On the other hand, the IPM orbit (2.2.1) was found to remain stable for all the values of $`N`$ and $`E`$ we studied (up to $`N=54`$ and $`E10^5`$). ## 3 Destabilization of SPOs and Globally Chaotic Dynamics Let us now study the chaotic behavior in the neighborhood of our unstable SPOs, starting with the well–known method of the evaluation of the spectrum of Lyapunov Exponents (LEs) of a Hamiltonian dynamical system, $`L_i,i=1,\mathrm{},2N`$ where $`L_1L_{\mathrm{max}}>L_2>\mathrm{}>L_{2N}`$. The LEs measure the rate of exponential divergence of initially nearby orbits in the phase space of the dynamical system as time approaches infinity. In Hamiltonian systems, the LEs come in pairs of opposite sign, so their sum vanishes, $`_{i=1}^{2N}L_i=0`$ and two of them are always equal to zero corresponding to deviations along the orbit under consideration. If at least one of them (the largest one) $`L_1L_{\mathrm{max}}>0`$, the orbit is chaotic, i.e. almost all nearby orbits diverge exponentially in time, while if $`L_{\mathrm{max}}=0`$ the orbit is stable (linear divergence of initially nearby orbits). Benettin et al., \[$`1980a,b`$\] studied in detail the problem of the computation of all LEs and proposed an efficient algorithm for their numerical computation, which we use here. In particular, $`L_iL_i(\stackrel{}{x}(t))`$ for a given orbit $`\stackrel{}{x}(t)`$ is computed as the limit for $`t\mathrm{}`$ of the quantities $`K_t^i`$ $`=`$ $`{\displaystyle \frac{1}{t}}\mathrm{ln}{\displaystyle \frac{\stackrel{}{w_i}(t)}{\stackrel{}{w_i}(0)}},`$ (48) $`L_i`$ $`=`$ $`\underset{t\mathrm{}}{lim}K_t^i`$ (49) where $`\stackrel{}{w_i}(0)`$ and $`\stackrel{}{w_i}(t),i=1,\mathrm{},2N`$ are deviation vectors from the given orbit $`\stackrel{}{x}(t)`$, at times $`t=0`$ and $`t>0`$ respectively. The time evolution of $`\stackrel{}{w_i}`$ is given by solving the so–called variational equations, i.e. the linearized equations about the orbit. Generally, for almost all choices of initial deviations $`\stackrel{}{w_i}(0)`$, the limit $`t\mathrm{}`$ of Eq. (49) gives the same $`L_i`$. In practice, of course, since the exponential growth of $`\stackrel{}{w_i}(t)`$ occurs for short time intervals, one stops the evolution of $`\stackrel{}{w_i}(t)`$ after some time $`T_1`$, records the computed $`K_{T_1}^i`$, orthogonormalizes the vectors $`\stackrel{}{w_i}(t)`$ and repeats the calculation for the next time interval $`T_2`$, etc. obtaining finally $`L_i`$ as an average over many $`T_j`$, $`j=1,2,\mathrm{},n`$ given by $$L_i=\frac{1}{n}\underset{j=1}{\overset{n}{}}K_{T_j}^i,n\mathrm{}.$$ (50) Next, we varied the values of the energy $`E`$ keeping $`N`$ fixed and studied the behavior of the Lyapunov exponents, using as initial conditions the OPMs (23) and (2.2.1) of the FPU and BEC Hamiltonians respectively. First, we find that the values of the maximum Lyapunov exponent $`L_1`$ increase by two distinct power–law behaviors ($`L_1E^B,B>0`$) as is clearly seen in Fig. 4 for the OPM of the FPU system. The result for the $`L_1`$ for the power–law behavior shown by solid line in Fig. 4 is in agreement with the results in Rechester et al., and Benettin, , where they obtain $`B=0.5`$ and $`B=2/3`$ for low dimensional systems and differs slightly from the one obtained in Livi et al., for the higher–dimensional case of $`N=80`$. We also find the same power–law behaviors with similar exponents $`B`$ for the other positive Lyapunov exponents as well. For example, for the $`L_2`$ we obtain $`L_2E^{0.62}`$ and $`L_2E^{0.48}`$ and $`L_3E^{0.68}`$ and $`L_3E^{0.49}`$, with the transition occurring at $`E28.21`$ for all of them. Turning now to the full Lyapunov spectrum in Fig. 5 we see that, for fixed $`N`$, as the energy is increased (and more eigenvalues of the monodromy matrix exit the unit circle) the Lyapunov spectrum tends to fall on a smooth curve for the OPM orbits of FPU and BEC, see Fig. 5(a), (b), as well as for the OHS mode with periodic boundary conditions (Fig. 5(c)). Observe that in Fig. 5(c) we have plotted the Lyapunov spectrum of both the OPM (23) of the FPU Hamiltonian (4) and of the OHS mode (6) for $`N=16`$ and periodic boundary conditions at the energy $`E=6.82`$ where both of them are destabilized and their distance in phase space is such that they are far away from each other. We clearly see that the two Lyapunov spectra are almost identical suggesting that their chaotic regions are somehow “connected”, as orbits starting initially in the vicinity of one of these SPOs visit often in the course of time the chaotic region of the other one. In Fig. 5(d) we have plotted the positive Lyapunov exponents spectra of three neighboring orbits of the OHS mode for $`N=15`$ dof in three different energies and observe that the curves are qualitatively the same as in Fig. 5(a) and (c). Finally, in Fig. 6 we have plotted the eigenvalues of the monodromy matrix of several SPOs and have observed the following: For the OPM of the FPU Hamiltonian the eigenvalues exit from $`1`$ (as the orbit destabilizes via period–doubling) and continue to move away from the unit circle, as $`E`$ increases further, see Fig. 6(a). By contrast, the eigenvalues of the OPM of the BEC Hamiltonian exit from $`+1`$ by a symmetry breaking bifurcation and for very large $`E`$ tend to return again to $`+1`$, see Fig. 6(b). This does not represent, however, a return to globally regular motion around this SPO, as the Lyapunov exponents in its neighborhood remain far from zero. Finally, in Fig. 6(c), we show an example of the fact that the eigenvalues of the IPM orbit of the BEC system, remain all on the unit circle, no matter how high the value of the energy is. Here, $`N=6`$, but a similar picture occurs for all the other values of $`N`$ we have studied up to $`N=54`$ and $`E10^5`$. ## 4 Using SALI to Estimate the “Size” of Islands of Regular Motion In this section, we estimate the “size” of islands of regular motion around stable SPOs using the Smaller Alignment Index (SALI) method Skokos, ; Skokos et al., 2003a ; Skokos et al., 2003b ; Skokos et al., , to distinguish between regular and chaotic orbits in the FPU and BEC Hamiltonians. The computation of the SALI has proved to be a very efficient method in revealing rapidly and with certainty the regular vs. chaotic nature of orbits, as it exhibits a completely different behavior for the two cases: It fluctuates around non–zero values for regular orbits, while it converges exponentially to zero for chaotic orbits. The behavior of the SALI for regular motion was studied and explained in detail by Skokos et al., 2003b , while a more analytical study of the behavior of the index in the case of chaotic motion can be found in Skokos et al., . As a first step, let us verify in the case of chaotic orbits of our $`N`$–degree of freedom systems, the validity of SALI’s dependence on the two largest Lyapunov exponents $`L_1L_{\mathrm{max}}`$ and $`L_2`$ proposed and numerically checked for $`N=2`$ and $`3`$, in Skokos et al., $$\mathrm{SALI}(\mathrm{t})e^{(L_1L_2)t}.$$ (51) This expression is very important as it implies that chaotic behavior can be decided by the exponential decay of this parameter, rather than the often questionable convergence of Lyapunov exponents to a positive value. To check the validity of (51) let us take as an example the OHS mode (6) of the FPU Hamiltonian (4) using fixed boundary conditions, with $`N=15`$ dof and $`\beta =1.04`$, at the energy $`E=21.6`$ and calculate the Lyapunov exponents, as well as the corresponding SALI evolution. Plotting SALI as a function of time $`t`$ (in linear scale) together with its analytical formula (51) in Fig. 7(a), we see indeed an excellent agreement. Increasing further the energy to the value $`E=26.6878`$, it is in fact possible to verify expression (51), even in the case where the two largest Lyapunov exponents are nearly equal, as Fig. 7(b) evidently shows! Exploiting now the different behavior of SALI for regular and chaotic orbits, we estimate approximately the “size” of regions of regular motion (or, “islands” of stability) in phase space, by computing SALI at points further and further away from a stable periodic orbit checking whether the orbits are still on a torus (SALI$`10^8`$) or have entered a chaotic “sea” (SALI$`<10^8`$) up to the integration time $`t=4000`$. The initial conditions are chosen perturbing all the positions of the stable SPO by the same quantity $`dq`$ and all the canonically conjugate momenta by the same $`dp`$ while keeping always constant the integral $`F`$, given by Eq. (33), in the case of the BEC Hamiltonian (32) and the energy $`E`$ in the case of the FPU Hamiltonian (4). In this way, we are able to estimate the approximate “magnitude” of the islands of stability for the OPM of Hamiltonian (4) and for the IPM and OPM of Hamiltonian (32) varying the energy $`E`$ and the number of dof $`N`$. In the case of the OPM solutions of both Hamiltonians, as the number of dof $`N`$ increases, for fixed energy $`E`$, the islands of stability eventually shrink to zero and the SPOs destabilize. For example, this is seen in Fig. 8(a) for the islands around the OPM (2.2.1) of the BEC Hamiltonian (32). A surprising behavior, however, is observed for the same SPO, if we keep $`N`$ fixed and increase the energy: Instead of diminishing, as expected from the FPU and other examples, the island of stability actually grows, as shown in Fig. 8(b), for the case of $`N=6`$ dof. In fact, it remains of considerable size until the SPO is destabilized for the first time, through period–doubling bifurcation at $`E3.1875`$, whereupon the island ceases to exist! But what happens to the island of stability around the IPM solution of the BEC Hamiltonian (32), which does not become unstable for all values of $`N`$ and $`E`$ we studied? Does it shrink to zero at sufficiently large $`E`$ or $`N`$? From Fig. 8(c) we see that for a fixed value of $`N`$, the size of this island also increases as the energy increases. In fact, this SPO has large islands about it even if the energy is increased, keeping the ratio $`E/N`$ constant (see Fig. 8(d)). This was actually found to be true for considerably larger $`E`$ and $`N`$ values than shown in this figure. ## 5 Lyapunov Spectra and the Thermodynamic Limit Finally, choosing again as initial conditions the unstable OPMs of both Hamiltonians, we determine some important statistical properties of the dynamics in the so–called thermodynamic limit of $`E`$ and $`N`$ growing indefinitely, while keeping $`E/N`$ constant. In particular, we compute the spectrum of the Lyapunov exponents of the FPU and BEC systems starting at the OPM solutions (23) and (2.2.1) for energies where these orbits are unstable. We thus find that the Lyapunov exponents are well approximated by smooth curves of the form $`L_iL_1e^{\alpha i/N}`$, for both systems, with $`\alpha 2.76`$, $`\alpha 3.33`$ respectively and $`i=1,2,\mathrm{},K(N)`$ where $`K(N)3N/4`$ (see Fig. 9). Specifically, in the case of the OPM (23) of the FPU Hamiltonian (4) for fixed energy density $`\frac{E}{N}=\frac{3}{4}`$ we find $$L_i(N)L_1(N)e^{2.76\frac{i}{N}},$$ (52) while, in the case of the OPM (2.2.1) of the BEC Hamiltonian (32) for fixed energy density $`\frac{E}{N}=\frac{3}{2}`$, a similar behavior is observed, $$L_i(N)L_1(N)e^{3.33\frac{i}{N}}.$$ (53) These exponential formulas were found to hold quite well, up to $`i=K(N)3N/4`$. For the remaining exponents, the spectrum is seen to obey different decay laws, which are not easy to determine. As this appears to be a subtle matter, however, we prefer to postpone it for a future publication. The functions (52) and (53), provide in fact, invariants of the dynamics, in the sense that, in the thermodynamic limit, we can use them to evaluate the average of the positive Lyapunov exponents (i.e. the Kolmogorov–Sinai entropy per particle) for each system and find that it is a constant characterized by the value of the maximum Lyapunov exponent $`L_1`$ and the exponent $`\alpha `$ appearing in them. In Fig. 10 we compute the well–known Kolmogorov–Sinai entropy $`h_{KS}(N)`$ Pesin, ; Hilborn, (solid curves), which is defined as the sum of the $`N1`$ positive Lyapunov exponents, $$h_{KS}(N)=\underset{i=1}{\overset{N1}{}}L_i(N),L_i(N)>0.$$ (54) In this way, we find, for both Hamiltonians, that $`h_{KS}(N)`$ is an extensive thermodynamic quantity as it is clearly seen to grow linearly with $`N`$ ($`h_{KS}(N)N`$), demonstrating that in their chaotic regions the FPU and BEC Hamiltonians behave as ergodic systems of statistical mechanics. Finally, using Eqs. (52) and (53), as if they were valid for all $`i=1,\mathrm{},N1`$, we approximate the sum of the positive Lyapunov exponents $`L_i`$, and calculate the $`h_{KS}(N)`$ entropy from Eq. (54) as $$h_{KS}(N)L_{\mathrm{max}}\frac{1}{1+e^{\frac{2.76}{N}}},L_{\mathrm{max}}0.095,(\mathrm{FPU}\mathrm{OPM})$$ (55) and $$h_{KS}(N)L_{\mathrm{max}}\frac{1}{1+e^{\frac{3.33}{N}}},L_{\mathrm{max}}0.34,(\mathrm{BEC}\mathrm{OPM}).$$ (56) In Fig. 10, we have plotted Eqs. (55) and (56) with dashed curves (adjusting the proportionality constants appropriately) and obtain nearly straight lines with the same slope as the data computed by the numerical evaluation of the $`h_{KS}(N)`$, from Eq. (54). ## 6 Conclusions In this paper we have investigated the connection between local and global dynamics of two $`N`$ dof Hamiltonian systems, describing $`1`$D nonlinear lattices with different origins known as the FPU and BEC systems. We focused on solutions located in the neighborhood of simple periodic orbits and showed that as the energy increases beyond the destabilization threshold, all positive Lyapunov exponents increase monotonically with two distinct power–law dependencies on the energy. We also computed the destabilization energy per particle of the OHS mode and of the two OPM orbits and found that it decays with a simple power–law of the form $`E_c/NN^\alpha `$, $`\alpha =`$ 1 or 2. One notable exception is the IPM orbit of the BEC Hamiltonian, which is found to be stable for any energy and number of dof we considered! Furthermore, we found that as we increase the energy $`E`$ of both Hamiltonians for fixed $`N`$, the SPOs behave in very different ways: In the OPM case of the FPU Hamiltonian the eigenvalues of the monodromy matrix always move away from $`1`$ with increasing energy, while in the OPM of the BEC Hamiltonian these eigenvalues return to $`+1`$, for very high energies. Nevertheless, this behavior does not represent a return of the system to a globally regular motion around the SPO as one might have thought. It simply reflects a local property of the SPO, as the orbits in its neighborhood still have positive Lyapunov exponents which are far from zero. We have also been able to estimate the “size” of the islands of stability around the SPOs with the help of a recently introduced very efficient indicator called SALI and have seen them to shrink, as expected, when increasing $`N`$ for fixed $`E`$ in the case of the OPM of the FPU Hamiltonian. Of course, when we continue increasing the energy, keeping the number of dof fixed, the OPMs destabilize and the islands of stability are destroyed. Unexpectedly, however, in the case of the OPM of the BEC Hamiltonian these islands were seen to grow in size if we keep $`N`$ fixed and increase $`E`$ up to the destabilization threshold, a result that clearly requires further investigation. The peculiarity of our BEC Hamiltonian is even more vividly manifested in the fact that the islands of stability about its IPM orbit never vanish, remaining actually of significant size, even for very high values of $`E`$ and $`N`$ (keeping $`E/N`$ fixed). Starting always near unstable SPOs, we also calculated the Lyapunov spectra characterizing chaotic dynamics in the regions of the OHS and OPM solutions of both Hamiltonians. Keeping $`N`$ fixed, we found energy values where these spectra were practically the same near SPOs which are far apart in phase space, indicating that the chaotic regions around these SPOs are visited ergodically by each other’s orbits. Finally, using an exponential law accurately describing these spectra, we were able to show, for both Hamiltonians, that the associated Kolmogorov–Sinai entropies per particle increase linearly with $`N`$ in the thermodynamic limit of $`E\mathrm{}`$ and $`N\mathrm{}`$ and fixed $`E/N`$ and, therefore, behave as extensive quantities of statistical mechanics. Our results suggest, however, that, even in that limit, there may well exist Hamiltonian systems with significantly sized islands of stability around stable SPOs (like the IPM of the BEC Hamiltonian), which must be excluded from a rigorous statistical description. It is possible, of course, that these islands are too small in comparison with the extent of the chaotic domain on a constant energy surface and their “measure” may indeed go to zero in the thermodynamic limit. But as long as $`E`$ and $`N`$ are finite, they will still be there, precluding the global definition of probability densities, ensemble averages and the validity of the ergodic hypothesis over all phase space. Clearly, therefore, their study is of great interest and their properties worth pursuing in Hamiltonian systems of interest to physical applications. ## 7 Acknowledgements This work was partially supported by the European Social Fund (ESF), Operational Program for Educational and Vocational Training II (EPEAEK II) and particularly the Program HERAKLEITOS, providing a Ph. D scholarship for one of us (C. A.). C. A. also acknowledges with gratitude the $`3`$ month hospitality, March–June $`2005`$, of the “Center for Nonlinear Phenomena and Complex Systems” of the University of Brussels. In particular, he thanks Professor G. Nicolis, Professor P. Gaspard and Dr. V. Basios for their instructive comments and useful remarks in many discussions explaining some fundamental concepts treated in this paper. The second author (T. B.) wishes to express his gratitude to the Max Planck Institute of the Physics of Complex Systems at Dresden, for its hospitality during his $`3`$ month visit March–June $`2005`$, when this work was completed. In particular, T. B. wants to thank Dr. Sergej Flach for numerous lively conversations and exciting arguments on the stability of multi–dimensional Hamiltonian systems. Useful discussions with Professors A. Politi, R. Livi, R. Dvorak, F. M. Izrailev and Dr. T. Kottos are also gratefully acknowledged. The third author (C. S.) was partially supported by the Research Committee of the Academy of Athens. ## 8 Figure Captions 1. The solid curve corresponds to the energy per particle $`\frac{E_c}{N}`$, for $`\beta =1.04`$, of the first destabilization of the OHS nonlinear mode (6) of the FPU system (4) obtained by the numerical evaluation of the Hill’s determinant in (19), while the dashed line corresponds to the function $`\frac{1}{N}`$. Note that both axes are logarithmic. 2. The solid curve corresponds to the energy per particle $`\frac{E_c}{N}`$, for $`\beta =1`$, of the first destabilization of the nonlinear OPM (23) of the FPU system (4) obtained by the numerical evaluation of the Hill’s determinant, while the dashed line corresponds to the function $`\frac{1}{N^2}`$. Note that both axes are logarithmic. 3. The solid curve corresponds to the energy per particle $`\frac{E_c}{N}`$ of the first destabilization of the OPM (2.2.1) of the BEC Hamiltonian (32) obtained by the numerical evaluation of the eigenvalues of the monodromy matrix of Eq. (2.2.2), while the dashed line corresponds to the function $`\frac{1}{N^2}`$. Note that both axes are logarithmic. 4. The two distinct power–law behaviors in the evolution of the maximum Lyapunov exponent $`L_1`$ as the energy grows for the OPM (23) of the FPU Hamiltonian (4) for $`N=10`$. A similar picture is obtained for the $`L_2`$ and $`L_3`$ also, with similar exponents and the transition occurring at the same energy value (see text). Note that both axes are logarithmic. 5. (a) The spectrum of the positive Lyapunov exponents for fixed $`N=10`$ of the OPM (23) of the FPU Hamiltonian (4) as the energy grows. (b) Also for $`N=10`$, the OPM (2.2.1) of the BEC Hamiltonian (32) yields a similar picture as the energy is increased. (c) The Lyapunov spectrum of the OPM (23) of the FPU Hamiltonian (4) for $`N=16`$ and the OHS mode (6) of the same Hamiltonian and $`N`$, for periodic boundary conditions practically coincide at $`E=6.82`$ where both of them are destabilized. (d) The Lyapunov spectrum of the FPU OHS mode (6) with fixed boundary conditions for $`N=15`$ as the energy grows presents as shape which is qualitatively similar to what was found for the SPOs of panel (c). 6. The eigenvalues $`\lambda _j,j=1,\mathrm{},2N`$ (a) of the OPM (23) of the FPU Hamiltonian (4) with $`N=10`$ dof. (b) The eigenvalues of the OPM (2.2.1) of the BEC Hamiltonian (32) with the same number of dof $`N`$. (c) The eigenvalues of the IPM (2.2.1) of the BEC Hamiltonian (32) with $`N=6`$ dof. 7. (a) The time evolution of the SALI (solid curve) and of the Eq. (51) (dashed line) at $`E=21.6`$ of the OHS mode (6) of the FPU Hamiltonian (4) with fixed boundary conditions. (b) Similar plot to panel (a) but for the larger energy $`E=26.6878`$ for which the two largest Lyapunov exponents $`L_1`$ and $`L_2`$ are almost equal while all the other positive ones are very close to zero. In both panels the agreement between the data (solid curve) and the derived function of Eq. (51) (dashed line) is remarkably good. Note that the horizontal axes in both panels are linear. 8. (a) “Size” of the islands of stability of the OPM (2.2.1) of the BEC Hamiltonian (32) for $`N=4,6`$ and $`8`$ dof and SPOs constant energy $`E=1`$ before the first destabilization (see Table 3). (b) “Size” of the islands of stability of the same Hamiltonian and SPO as in (a) for $`N=6`$ dof and three different energies of the SPO before the first destabilization (see Table 3). (c) “Size” of the islands of stability of the same Hamiltonian as in (a) of the IPM (2.2.1) for $`N=6`$ dof and four different energies of the SPO. (d) “Size” of the islands of stability of the same Hamiltonian as in (a) of the IPM (2.2.1) for $`\frac{E}{N}=\frac{10}{3}`$. Here $`E`$ corresponds to the energy of the IPM. 9. (a) Positive Lyapunov exponents spectrum of the OPM (23) of the FPU Hamiltonian (4) for fixed $`\frac{E}{N}=\frac{3}{4}`$. (b) Positive Lyapunov exponents spectrum of the OPM (2.2.1) of the BEC Hamiltonian (32) for fixed $`\frac{E}{N}=\frac{3}{2}`$. In both panels $`i`$ runs from $`1`$ to $`N`$. 10. (a) The $`h_{KS}(N)`$ entropy of the OPM (23) of the FPU Hamiltonian (4) for fixed $`\frac{E}{N}=\frac{3}{4}`$ (solid curve) and the approximated formula (55) (dashed curve). (b) The $`h_{KS}(N)`$ entropy of the OPM (2.2.1) of the BEC Hamiltonian (32) for fixed $`\frac{E}{N}=\frac{3}{2}`$ (solid curve) and the approximated formula (56) (dashed curve).
warning/0506/quant-ph0506048.html	ar5iv	text	# Evanescence in Coined Quantum Walks ## 1 Introduction The first authors to discuss the quantum walk were Aharonov, Davidovich and Zagury in where they described a very simple realization in quantum optics. In this model a particle takes unit steps on the integers at each time step, starting at the origin. In Meyer proved that an additional spin-like degree of freedom was essential if the behaviour of the system was to be both unitary and non-trivial. Without this degree of freedom, the only way the evolution of the walk can avoid being purely ballistic is to relax the unitarity condition. This spin-like degree of freedom is sometimes called the chirality, or the coin, which is why this type of walk is sometimes called a “coined” walk . This is in sharp contrast to the continuous time walk which does not need a coin and which we will not discuss here. The chirality can take the values RIGHT and LEFT, or a coherent superposition of these. For a detailed introduction to quantum walks, we refer the reader to the review article in . Meyer and subsequent authors have considered two approaches to the discrete-time quantum walk, the path-integral approach of Feynman and the Schrödinger wave-mechanics approach, which reflect two complementary ways of formulating quantum mechanics . We refer to the paper by Ambainis, Bach, Nayak, Vishwanath and Watrous for proper definitions and more references. Both approaches are discussed in the paper of Ambainis et al. The probability distribution for this walk is shown in Figure 1, after enough time has elapsed for the asymptotic behaviour to manifest. This paper began as a sequel to the work of Carteret, Ismail and Richmond concerning the one-dimensional quantum walk on the integers, and contains the completion of the analysis of this quantum walk for the remaining exponential decay region in the Schrödinger picture. We believe the analysis presented in this paper to be interesting for three reasons. One is methodological; while analysing this system we encountered various links between a number of different methods in combinatorics which do not seem to be widely known, and which may be of use to the quantum information community when analysing more complicated systems than the one discussed here. The second motivation is rather more abstract. It is one of the fundamental principles of quantum mechanics that the wave-mechanics and path-integral representations of a system should produce exactly the same results. The quantum walk has been proposed as the quantum analogue of the classical random walk, in the hopes of ultimately defining a systematic procedure for “quantizing” classical random walk algorithms . Quantizing classical systems is something that must be done with considerable care; the obvious approach isn’t necessarily the correct one. While it is true that when such pathologies have been discovered in the past, they were found in much more exotic systems than this one, it is as well to check. One way to perform such a check is to verify that the Feynman equivalence principle still holds between the path-integral and wave-mechanics representations. As it happens, the results obtained from the two approaches do not at first appear to be equivalent. In fact they are, though the proof of their equivalence is nontrivial, and will be given in Section 3. In the course of this analysis we uncovered a small, but potentially significant omission in previous analyses of this system, which we will describe below. Another reason for performing this check is to explore the little mystery left at the end of . While the results from the path integral calculation made intuitive sense, the partial results from the wave-mechanics calculation for the exponential decay region were rather unexpected. Specifically, we found what appeared to be evanescent waves in this exponential decay region, which seemed to imply the presence of some kind of absorption mechanism, despite the fact that the definition of the model precludes any barrier or other source of dissipation in the system to cause these by the familiar mechanisms. So, we will check that the Feynman equivalence holds to verify that those evanescent waves are not simply some kind of mathematical artefact. We would also like to gain some insight into the physical interpretation of the mathematical behaviour of this part of the wave-function. A link between the behaviour of the quantum walk on the line and certain phenomena in quantum optics has been suggested previously by Knight, Roldan and Sipe in , and refined by Kendon and Sanders in . This connection will be discussed in more detail in Section 4. We will also describe a potentially useful method for obtaining integral representations of orthogonal polynomials from their generating functions using Lagrange Inversion. This bypasses the need to use the Darboux method and makes it possible to obtain uniformly convergent asymptotics directly from the generating function. We have included this in the Appendices in Subsection 6.2. ### 1.1 Some previous results In this section we will mention some results by other authors which we will have occasion to use later in this paper. One of the reasons for doing this is that different authors have used different labelling conventions and this will enable us to establish a consistent notation for use when we combine results from different papers with mutually incompatible conventions. We will state our results using the conventions in . Several early results in the theory of quantum walks are due to Meyer , who considered the wavefunction as a two-component vector of amplitudes of the particle being at point $`n`$ at time $`t`$. Let $$\mathrm{\Psi }(n,t)=\left(\begin{array}{c}\psi _R(n,t)\\ \psi _L(n,t)\end{array}\right)$$ (1) where the chirality of the top component is labelled RIGHT and the bottom LEFT. At each step the chirality of the particle evolves according to a unitary Hadamard transformation $`\|R`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\|R+\|L\right)`$ (2) $`\|L`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\|R\|L\right),`$ (3) which is why this quantum walk is sometimes called the “Hadamard” walk. The particle (or “walker”) then moves according to its new chirality state. Therefore, the particle obeys the recursion relations $`\mathrm{\Psi }_R(n,t+1)`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\mathrm{\Psi }_L(n+1,t)+{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{\Psi }_R(n1,t)`$ (4) $`\mathrm{\Psi }_L(n,t+1)`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\mathrm{\Psi }_L(n+1,t)+{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{\Psi }_R(n1,t).`$ (5) Meyer approached this problem from the path-integral point of view, and obtained expressions for the $`\psi `$-functions in terms of Jacobi polynomials. The standard notation in for Jacobi polynomials is $`P_n^{(\alpha ,\beta )}(z)`$ but we wish to follow the conventions in and subsequent papers that have now become standard in the literature on quantum walks, and use $`\alpha =n/t.`$ We will therefore use the notation $`J_q^{(r,s)}(w).`$ We find, in and when $`nt`$, ###### Theorem 1 (Ambainis et al. , after Meyer, ). $$\psi _R(n,t)(1)^{(tn)/2}=\{\begin{array}{cc}2^{n/21}J_{((t+n)/21)}^{(1,n)}(0)\hfill & \text{when }tn<0\hfill \\ \left(\frac{t+n}{tn}\right)2^{n/21}J_{(tn)/21}^{(1,n)}(0)\hfill & \text{when }0n<t\hfill \end{array}$$ (6) Also $$\psi _L(n,t)(1)^{(tn)/2}=\{\begin{array}{cc}2^{n/2}J_{(t+n)/2}^{(0,n1)}(0)\hfill & \text{when }tn<0\hfill \\ 2^{n/21}J_{(tn)/21}^{(0,n+1)}(0)\hfill & \text{when }0n<t\hfill \end{array}$$ (7) and $`\|\psi _R(n,t)\|^2`$ $`=\left({\displaystyle \frac{tn}{t+n}}\right)^2\|\psi _R(n,t)\|^2`$ (8) $`\|\psi _L(n,t)\|^2`$ $`=\|\psi _L(2n,t)\|^2.`$ (9) Ambainis et al. use the other sign convention, so one should interchange $`L`$ and $`R`$ (or equivalently, replace $`n`$ by $`n`$) to reflect the walk before comparing their results with ours. This is just a relabelling, and so their results can be stated as in the following theorem. We will prove the above results in the form below ###### Theorem 2 (Ambainis et al. ). When $`nt(\mathrm{mod}\mathrm{\hspace{0.17em}2})`$ and $`J_q^{(r,s)}(w)`$ denotes a Jacobi polynomial, then $$\psi _R(n,t)(1)^{(tn)/2}=\{\begin{array}{cc}(1)^{n+1}2^{n/2}J_{(tn)/2}^{(0,n1)}(0)\hfill & \text{when }0nt\hfill \\ (1)^{n+1}2^{n/21}J_{(t+n)/21}^{(0,n+1)}(0)\hfill & \text{when }t<n<0\hfill \end{array}$$ (10) Also $$\psi _L(n,t)(1)^{(tn)/2}=\{\begin{array}{cc}2^{n/21}J_{(tn)/21}^{(1,n)}(0)\hfill & \text{when }0n<t\hfill \\ \left(\frac{tn}{t+n}\right)2^{n/21}J_{(t+n)/21}^{(1,n)}(0)\hfill & \text{when }t<n<0.\hfill \end{array}$$ (11) and $`\psi _R(n,t)`$ $`=(1)^{n+1}\psi _R(n+2,t),`$ (12) $`(tn)\psi _L(n,t)`$ $`=(1)^n(t+n)\psi _L(n,t).`$ (13) A few Remarks: 1. Note that Theorem 2 differs from Theorem 1 by an external phase which has been dropped in previous analyses of this system; we state the symmetry relations for the $`\psi `$-functions rather than for their moduli-squared, as in . We will discuss this in more detail below, as some properties of Jacobi polynomials are required for the calculation. 2. There is a sign error in the symmetry relations for $`\psi _R`$ and $`\psi _L(n,t)`$ in Carteret et al. (which has been corrected in the arxiv version ). The symmetry relations will be proved in Lemmas 2, 3 and equation (77) of Subsection 3.1 using some integral representations of $`\psi _R(n,t)`$ and $`\psi _L(n,t)`$. 3. The endpoints where $`n=\pm t`$ have to be handled separately, see . For the starting conditions $`\psi _L(0,0)=1,`$ $`\psi _R(0,0)=0`$ the wavefunctions at the end-points (where $`n=\pm t`$) are $`\psi _R(t,t)`$ $`=(1)^{t+1}2^{t/2}`$ $`t=0,1,2,\mathrm{},`$ (14) $`\psi _L(t,t)`$ $`=0,`$ $`t=1,2,3,\mathrm{},`$ (15) $`\psi _R(t,t)`$ $`=0,`$ $`t=1,2,3,\mathrm{},`$ (16) $`\psi _L(t,t)`$ $`=(1)^t2^{t/2},`$ $`t=0,1,2,\mathrm{}.`$ (17) 4. The two different cases in (11) and (10) for $`n0`$ and $`n<0`$ can be combined into one case for all $`n`$ satisfying $`tn<t.`$ We prove this later using a symmetry property of the Jacobi polynomials. Our results in equations (12) and (13) are refinements of the corresponding relations in (9) and (8), after performing the relabelling necessary to compare results with different sign conventions. The asymptotic behaviour for the path-integral representation has been determined in Carteret et al. , starting from Theorem 2. The steepest descent technique was used on the standard integral representation for the Jacobi polynomial. The result was uniform exact asymptotics $`\alpha `$ in the range $`\|\alpha \|<1\epsilon ,`$ where $`\epsilon `$ is any positive number, in terms of Airy functions. This technique was used earlier for Jacobi polynomials by Saff and Varga and by Gawronkski and Sawyer ; however the connection with Airy functions had not been recognized as far as we know. The Airy function description is useful for $`\|\alpha \|`$ near $`1/\sqrt{2}`$ where the asymptotic behaviour changes from an oscillating cosine term times $`t^{1/2}`$ (for $`\|\alpha \|<1/\sqrt{2}`$) to exponentially small (for $`2^{1/2}+\epsilon <\|\alpha \|<1\epsilon ).`$ In this paper we analyze the Hadamard walk from the Schrödinger wave-mechanics point of view. The earliest work on this that the authors are aware of is that by Nayak and Vishwanath . They define $$\stackrel{~}{\mathrm{\Psi }}(\theta ,t)=\underset{n}{}\psi (n,t)e^{i\theta n},$$ (18) (where we have used the symbol $`\theta `$ for the momentum instead of the $`k`$ used in ) so the recursion relations above becomes $$\stackrel{~}{\mathrm{\Psi }}(\theta ,t+1)=M_\theta \stackrel{~}{\mathrm{\Psi }}(\theta ,t),$$ (19) where $$M_\theta =\left(\begin{array}{cc}e^{i\theta }& e^{i\theta }\\ e^{i\theta }& e^{i\theta }\end{array}\right).$$ (20) Thus $$\stackrel{~}{\mathrm{\Psi }}(\theta ,t)=M_\theta ^t\stackrel{~}{\mathrm{\Psi }}(\theta ,0),\stackrel{~}{\mathrm{\Psi }}(\theta ,0)=(1,0)^\mathrm{T},$$ (21) where the symbol $`^\mathrm{T}`$ denotes transposition. They show that the eigenvalues of $`M_\theta `$ are $`e^{i\omega _\theta }`$ and $`e^{i\omega _\theta }`$ where $`\omega _\theta `$ is the angle in $`[\pi /2,\pi /2]`$ such that $`\mathrm{sin}(\omega _\theta )=(\mathrm{sin}\theta )/\sqrt{2}.`$ They also use the other sign convention, so one should relabel $`L`$ and $`R`$ as before. Their results can be stated as in the following theorem. ###### Theorem 3 (Nayak and Vishwanath ). Let $`\alpha =n/t`$. Then $$\psi _R(n,t)=\frac{1+(1)^{n+t}}{2}\frac{1}{2\pi }_\pi ^\pi \frac{e^{i\theta }}{\sqrt{1+\mathrm{cos}^2\theta }}e^{i(\omega _\theta +\theta \alpha )t}𝑑\theta .$$ (22) $$\psi _L(n,t)=\frac{1+(1)^{n+t}}{2}\frac{1}{2\pi }_\pi ^\pi \left(1+\frac{\mathrm{cos}\theta }{\sqrt{1+\mathrm{cos}^2\theta }}\right)e^{i(\omega _\theta +\theta \alpha )t}𝑑\theta ,$$ (23) We will derive the asymptotic behaviour of the $`\psi `$-functions starting from Theorem 3 in Section 2. We will only give the complete details for the exponential decay range, as the calculation for the oscillatory region has already been done by others . The conventional version of the method of stationary phase, as used by Nayak-Vishwanath , does not work in the exponentially small region, as the stationary points of the phase function have left the real line. We will show how to extend and refine this technique so that it can be made to work in this situation; the modification is an application of the method of steepest descents. It is not obvious that the formulæ obtained by each method are the same, but we will prove below in Section 3 that they are. This means that precisely the same asymptotic behaviour can be found using both the path-integral and wave-mechanics descriptions of quantum mechanics. ## 2 Generalizing the method of stationary phase The aim of this section is to extend the method of stationary phase so that it can cope with stationary points that occur as complex conjugate pairs on either side of the real axis. We start with the integral representation of $`\psi _R`$ in Theorem 3. In this representation the integration is performed along the real axis. Nayak and Vishwanath consider the case $`\|\alpha \|<1/\sqrt{2}`$ when there are two stationary points (defined below) inside the interval of integration $`[\pi ,\pi ]`$. When we find the critical points of the phase function, we obtain an equation for $`\theta `$ (called $`k`$ in ) at the critical points as a function of $`\alpha ,`$ which is $$\mathrm{cos}\theta _\alpha =\frac{\alpha }{\sqrt{1\alpha ^2}}$$ (24) from which the critical value of $`\omega _\theta ,`$ call it $`\omega _{\theta _\alpha },`$ can be obtained using the Pythagoras rule ($`\mathrm{cos}^2\theta +\mathrm{sin}^2\theta =1`$) and $`\omega _\theta =\mathrm{arcsin}\frac{\mathrm{sin}\theta }{\sqrt{2}}`$ from . However, when $`\|\alpha \|>1/\sqrt{2}`$ this equation no longer has any real solutions and thus the corresponding stationary points are no longer on the real axis, see figure 2. Therefore the standard method of stationary phase cannot provide the exact asymptotics. The stationary points have “moved” off the real axis and become a complex conjugate pair. We shall move the contour of integration to follow the stationary points, whilst ensuring that the contour still goes through one of them. Note also that the stationary points become saddle-points on leaving the real line; we will return to this fact in Subsection 2.2, below. ### 2.1 The function $`𝝎_𝜽`$ in the complex plane The key to evaluating the asymptotics for these integrals lies in the behaviour of the phase function $`\omega _\theta .`$ We will therefore begin by describing the analytic properties of the function $`\omega _\theta `$ in the complex plane, in particular in the strip $`\pi \mathrm{}\theta \pi `$. We will need this information when we replace the initial interval of integration by a contour in the complex plane, as explained in the next subsection. We will also need estimates of $`\omega _\theta `$ at $`+\mathrm{}`$ in this strip to show that our new contour integral converges. The singular points of $`\omega _\theta =\mathrm{arcsin}\left(\frac{\mathrm{sin}\theta }{\sqrt{2}}\right)`$ are found from the equations $`\frac{\mathrm{sin}\theta }{\sqrt{2}}=\pm 1`$. When we write $`\theta =u+iv`$, with $`u[\pi ,\pi ]`$ and $`v,`$ we conclude from the equations $`\mathrm{sin}(u+iv)=\pm \sqrt{2},`$ where $$\mathrm{sin}(u+iv)=\mathrm{sin}u\mathrm{cosh}v+i\mathrm{cos}u\mathrm{sinh}v$$ (25) that $`u=\pm \frac{1}{2}\pi `$ and $`\mathrm{cosh}v=\sqrt{2}`$ (or $`v=\pm \mathrm{arcsinh}(1)`$). Because of conjugation and symmetry there are four singular points, namely $`\pm \frac{1}{2}\pi +i\mathrm{arcsinh}(1)`$ and $`\pm \frac{1}{2}\pi i\mathrm{arcsinh}(1)`$. On the four half lines $`\pm \frac{1}{2}\pi +iv`$ with $`\|v\|>\mathrm{arcsinh}(1)`$ the function $`\mathrm{sin}(u+iv)`$ is real, and has an absolute value greater than $`\sqrt{2}`$. These four half-lines are taken as branch cuts of the multi-valued function $`\omega _\theta `$. They correspond with the two branch cuts of the function $`\mathrm{arcsin}z=\mathrm{arcsin}(x+iy)`$ in the $`z`$plane from $`x=\pm 1`$ to $`x=\pm \mathrm{}`$, with $`y=0`$. The strip $`\pi u\pi `$ that is delineated by these branch cuts is the principal Riemann sheet on which $`\omega _\theta `$ is analytic and single-valued. We consider the principal branch of this function that is real on $`[\pi ,\pi ]`$ and continuously extended on the principal sheet. Since the function is periodic, the same holds for the other strips $`[k\pi ,(k+2)\pi ]`$, $`k\text{Z}\text{Z}`$. In Figure 3 we show the conformal mapping by $`\omega _\theta `$ from the strip $`\pi u\pi `$. We show the images of a number of lines, where we concentrate on $`v0`$. For $`v0`$ a similar picture can be given. We observe the following useful facts. 1. The image of the interval $`[0,\pi ]`$ must go around a branch cut because $`\omega _\theta `$ is not single valued on this interval; the image point $`A`$ is given by $`A=\mathrm{arcsin}\frac{1}{\sqrt{2}}`$. 2. The points $`D`$ and $`F`$ are on different sides of the branch cut; the loop $`DEF`$ around the branch cut is mapped to the vertical $`DEF,`$ and the same goes for the three points $`D^{}E^{}F^{}`$. 3. On the positive imaginary axis $`u=0,v0`$ $`\omega _\theta `$ has the form (cf. (25)) $$\omega _\theta =i\mathrm{arcsinh}\frac{\mathrm{sinh}v}{\sqrt{2}}=i\mathrm{ln}\left(\frac{\mathrm{sinh}v}{\sqrt{2}}+\sqrt{\frac{\mathrm{sinh}^2v}{2}+1}\right),v0.$$ (26) 4. On the half-lines $`u=\pm \pi ,v0,\omega _\theta `$ has the form $$\omega _\theta =i\mathrm{arcsinh}\frac{\mathrm{sinh}v}{\sqrt{2}}=i\mathrm{ln}\left(\frac{\mathrm{sinh}v}{\sqrt{2}}+\sqrt{\frac{\mathrm{sinh}^2v}{2}+1}\right),v0.$$ (27) 5. For $`\theta `$ on the vertical $`ED^{}`$ we can write $`\omega _\theta `$ in the form $$\omega _\theta =\frac{1}{2}\pi +i\mathrm{ln}\left(\frac{\mathrm{cosh}v}{\sqrt{2}}+\sqrt{\frac{\mathrm{cosh}^2v}{2}1}\right),\mathrm{cosh}v\sqrt{2}.$$ (28) 6. For $`\theta `$ on $`ED`$ we must choose the negative square root, which gives $$\omega _\theta =\frac{1}{2}\pi i\mathrm{ln}\left(\frac{\mathrm{cosh}v}{\sqrt{2}}+\sqrt{\frac{\mathrm{cosh}^2v}{2}1}\right),\mathrm{cosh}v\sqrt{2}.$$ (29) We will also need the following lemma, in order to bound $`\omega _\theta `$ as the imaginary part of $`\theta `$ tends to $`+\mathrm{},`$ and hence guarantee that these integrals will converge. ###### Lemma 1. If $`\theta =u+iv,v>0`$, then, as $`v+\mathrm{}`$ $$e^{i\omega _\theta t}=\{\begin{array}{cc}𝒪(e^{+vt}),\hfill & \|u\|<\frac{1}{2}\pi ,\hfill \\ 𝒪(e^{vt}),\hfill & \frac{1}{2}\pi <\|u\|\pi .\hfill \end{array}$$ (30) ###### Proof. Since $`\omega _\theta =\mathrm{arcsin}\left(\frac{\mathrm{sin}\theta }{\sqrt{2}}\right)`$ we have $`\mathrm{sin}\omega _\theta =\mathrm{sin}\theta /\sqrt{2}`$; so $$e^{i\omega _\theta }=\mathrm{cos}\omega _\theta i\frac{\mathrm{sin}\theta }{\sqrt{2}}.$$ (31) Hence $$e^{i\omega _\theta }=\pm \sqrt{1\frac{\mathrm{sin}^2\theta }{2}}i\frac{\mathrm{sin}\theta }{\sqrt{2}}$$ (32) where the $`\pm `$ sign in front of the first term has yet to be determined. For small values of $`\theta `$, it is obvious that we should select the $`+`$ sign, because both sides of equation (32) have to approach unity as $`\theta 0`$. In fact, the $`+`$ sign should be chosen throughout the strip $`\|u\|<\frac{1}{2}\pi `$. This is because $$1\frac{\mathrm{sin}^2\theta }{2}=\frac{3+\mathrm{cos}2\theta }{4}\frac{1}{8}e^{2i\theta }$$ (33) as $`\theta +i\mathrm{}`$, we conclude that in the strip $`\|u\|<\frac{1}{2}\pi `$ $$e^{i\omega _\theta }\frac{1}{2\sqrt{2}}e^{i\theta }i\frac{\mathrm{sin}\theta }{\sqrt{2}}\frac{1}{\sqrt{2}}e^{i\theta }=𝒪(e^v),$$ (34) as $`v+\mathrm{}`$. Observe that this is in agreement with the behaviour of $`\omega _\theta `$ on the positive imaginary axis, as given in equation (26). It also agrees with the conformal mapping shown in Figure 3, where we see that the domain $`AEFGF^{}E^{}A^{}OA`$ is mapped to $`\mathrm{}\omega _\theta >0`$. The figure also shows that the domains $`ABCDEA`$ and $`A^{}B^{}C^{}D^{}E^{}A^{}`$ are mapped to $`\mathrm{}\omega _\theta <0`$. This corresponds to choosing the negative values for $`\theta `$ in (32) in these domains; this gives the estimate in the second line of (30). This proves the lemma. ∎ Now that we have established the behaviour of the phase function $`\omega _\theta ,`$ we can proceed to choose an appropriate contour of integration. ### 2.2 Saddle-point analysis We will now obtain an asymptotic approximation for $`\psi _R`$ in the exponentially small range outside the main peaks. To find convenient locations for the contours of integration with respect to the stationary points, we will begin with the integral in (22) of Theorem 3 for $`\psi _R(2n,t)`$, and use the symmetry rule for $`\psi _R(n,t)`$ (cf. Theorem 2) to obtain the result for $`\psi _R(n,t)`$. So, our starting point is $$\psi _R(n,t)=(1)^{n+1}\psi _R(2n,t)=\frac{(1)^{n+1}}{2\pi }_\pi ^\pi \frac{e^{i\theta }}{\sqrt{1+\mathrm{cos}^2\theta }}e^{i(\omega _\theta \theta \alpha )t}𝑑\theta ,$$ (35) where $`\alpha =n/t`$. We have dropped the factor $`\frac{1+(1)^{n+t}}{2}`$, because we always can assume that $`n`$ and $`t`$ have the same parity; the wavefunction is identically zero otherwise as the walker must always move at each time-step. We first locate the stationary points or saddle-points in the traditional way, that is, we solve the equation $$\alpha =\frac{d\omega _\theta }{d\theta }=\frac{(\mathrm{cos}\theta _\alpha )/\sqrt{2}}{\sqrt{1\mathrm{sin}^2\theta _\alpha /2}}=\frac{1}{\sqrt{1+\mathrm{cos}^2\theta _\alpha }}\mathrm{cos}\theta _\alpha .$$ (36) Note that in (36) $`\mathrm{cos}\theta _\alpha `$ and $`\alpha `$ have the same sign (and $`\alpha `$ is positive). This gives $$\mathrm{cos}\theta _\alpha =\pm \frac{\alpha }{\sqrt{1\alpha ^2}}.$$ (37) Thus, when $`\alpha <1/\sqrt{2},`$ this gives two real stationary points $$\theta _\alpha =\pm \mathrm{arccos}\left(\frac{\alpha }{\sqrt{1\alpha ^2}}\right),$$ (38) which are used in the stationary phase method in . If $`1/\sqrt{2}<\alpha <1`$ these points are purely imaginary, and they are given by $$\theta _\alpha =\pm i\text{arccosh}\left(\frac{\alpha }{\sqrt{1\alpha ^2}}\right).$$ (39) When $`1/\sqrt{2}<\|\alpha \|<1`$ we shift the contour in the integral representation of the $`\psi _R`$ given in (35) off the real axis to the segments shown in Figure 5. Our modified stationary phase method is in fact a version of the method of steepest descents. The contour of integration goes through the saddle-point on the positive imaginary axis, that is, through $`\theta _\alpha =i\mathrm{arccosh}(\alpha /\sqrt{1\alpha ^2})`$ and fixes the imaginary part of $`i\omega _\theta i\theta \alpha .`$ This is equivalent to fixing the real part of $`\omega _\theta \theta \alpha .`$ We proceed as follows. Consider the integral along the contour in Figure 5. We can make this into a closed contour by adding in segments from $`\pi \pi ,`$ $`\pi \pi +i\mathrm{}`$ and $`\pi +i\mathrm{}\pi `$, thus obtaining an integral over $$(\pi ,\pi )(\pi ,\pi +i\mathrm{})(\pi +i\mathrm{},\theta _\alpha )(+\theta _\alpha ,\pi +i\mathrm{})(\pi +i\mathrm{},\pi ).$$ (40) The singular points of $`\sqrt{1+\mathrm{cos}^2\theta }`$ follow from solving $`\mathrm{cos}^2\theta =1`$, which gives $`\theta =\pm \pi /2\pm i\mathrm{arcsinh}(1)`$ (see the open dots in Figure 5). We avoid the singularities and branch cuts of the square root in the integrand and of the function $`\omega _\theta `$ (these singularities are the same for both functions; see Figure 3). The integrand is then analytic around and inside the contour in (40), so the integral around the contour is zero. The integrals over the curves indicated below are therefore equal, $$(\pi ,\pi )=(\pi ,\pi +i\mathrm{})(\pi +i\mathrm{},\theta _\alpha )(+\theta _\alpha ,\pi +i\mathrm{})(\pi +i\mathrm{},\pi ).$$ (41) The steepest-descent curves for this integral are shown in Figure 5. Furthermore the periodicity of the integrand modulo $`2\pi `$ means that the segment integrals from $`\pi `$ to $`\pi +i\mathrm{}`$ and from $`\pi +i\mathrm{}`$ to $`\pi `$ cancel, so these are not shown in Figure 5. Thus, the integral from $`\pi `$ to $`\pi `$ equals the integral along the contour from $`\pi +i\mathrm{}`$ to $`\pi +i\mathrm{}`$ through $`\theta _\alpha .`$ From Lemma 1 we conclude that $$e^{i(\omega _\theta \alpha )t}=𝒪\left(e^{(1\alpha )vt}\right)$$ (42) as $`v+\mathrm{}`$ in the strips $`\pi u<\frac{1}{2}\pi `$ and $`\frac{1}{2}\pi <u\pi `$. Hence, convergence at infinity on a contour as shown in Figure 5 is guaranteed. ### 2.3 Evaluating the main contribution We evaluate $`e^{i\omega _\alpha +i\theta _\alpha \alpha }`$ for the saddle-point on the positive imaginary axis, that is, for $$\theta _\alpha =i\mathrm{arccosh}\left(\frac{\alpha }{\sqrt{1\alpha ^2}}\right).$$ (43) Now, with $`x=\alpha /\sqrt{1\alpha ^2}`$, we obtain $$\theta _\alpha =i\mathrm{arccosh}x=i\mathrm{ln}(x+(x^21)^{1/2})=i\mathrm{ln}\left(\frac{\alpha +\sqrt{2\alpha ^21}}{\sqrt{1\alpha ^2}}\right).$$ (44) Thus $$e^{i\theta _\alpha \alpha }=\mathrm{exp}\left(\alpha \mathrm{ln}\left(\frac{\alpha +\sqrt{2\alpha ^21}}{\sqrt{1\alpha ^2}}\right)\right)=\left(\frac{\sqrt{1\alpha ^2}}{\alpha +\sqrt{2\alpha ^21}}\right)^\alpha .$$ (45) Let us now consider $`\omega _\alpha =\mathrm{arcsin}\left(\frac{\mathrm{sin}\theta _\alpha }{\sqrt{2}}\right)`$. We have $$\mathrm{cos}\theta _\alpha =\frac{\alpha }{\sqrt{1\alpha ^2}},$$ (46) so therefore $$\mathrm{sin}\theta _\alpha =\sqrt{1\alpha ^2/(1\alpha ^2)}=\sqrt{\frac{12\alpha ^2}{1\alpha ^2}}=i\sqrt{\frac{2\alpha ^21}{1\alpha ^2}}.$$ (47) Thus we obtain $$\omega _\alpha =i\mathrm{arcsinh}\left(\frac{1}{\sqrt{2}}\sqrt{\frac{2\alpha ^21}{1\alpha ^2}}\right).$$ (48) Now $`\mathrm{arcsinh}x=\mathrm{ln}(x+\sqrt{x^2+1})`$, so $`\omega _\alpha `$ $`=i\mathrm{ln}\left({\displaystyle \frac{\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}}+\left({\displaystyle \frac{2\alpha ^21}{2(1\alpha ^2)}}+1\right)^{1/2}\right)`$ $`=i\mathrm{ln}\left({\displaystyle \frac{\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}}+{\displaystyle \frac{(2\alpha ^21+22\alpha ^2)^{1/2}}{\sqrt{2(1\alpha ^2)}}}\right)`$ $`=i\mathrm{ln}\left({\displaystyle \frac{1+\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}}\right).`$ (49) Thus $$e^{i\omega _\alpha }=\mathrm{exp}\left(\mathrm{ln}\frac{1+\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}\right)=\frac{1+\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}.$$ (50) and hence $$e^{i\omega _\alpha +i\theta _\alpha \alpha }=\left(\frac{\sqrt{1\alpha ^2}}{\alpha +\sqrt{2\alpha ^21}}\right)^\alpha \frac{1+\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}.$$ (51) Since $`\omega _\theta \theta _\alpha \alpha `$ is an odd function of $`\theta _\alpha `$, we will obtain the reciprocal of this result for the saddle-point $`i\theta _\alpha .`$ The saddle-point on the positive imaginary axis (when $`\alpha >1/\sqrt{2}`$) is indeed the relevant one as we will see. The exact saddle-point contours are shown in Figures 4 and 5. Now that we have established these preliminary results, we can proceed to prove ###### Theorem 4. If $`1/\sqrt{2}+\epsilon <\|\alpha \|<1\epsilon ,`$ then $$\psi _R(n,t)\frac{(1)^{n+1}\left(\alpha +\sqrt{2\alpha ^21}\right)t^{1/2}}{\sqrt{2\pi (1\alpha ^2)\sqrt{2\alpha ^21}}}\left(\left(\frac{\sqrt{1\alpha ^2}}{\alpha +\sqrt{2\alpha ^21}}\right)^\alpha \frac{1+\sqrt{2\alpha ^21}}{2\sqrt{1\alpha ^2}}\right)^t.$$ (52) $$\psi _L(n,t)\frac{(1)^n(1\alpha )t^{1/2}}{\sqrt{2\pi (1\alpha ^2)\sqrt{2\alpha ^21}}}\left(\left(\frac{\sqrt{1\alpha ^2}}{\alpha +\sqrt{2\alpha ^21}}\right)^\alpha \frac{1+\sqrt{2\alpha ^21}}{2\sqrt{1\alpha ^2}}\right)^t.$$ (53) ###### Proof. We prove the equation for $`\psi _R,`$ the proof of the equation for $`\psi _L`$ is very similar. Since $`\omega _\theta =\mathrm{arcsin}\left(\frac{\mathrm{sin}\theta }{\sqrt{2}}\right),`$ we have $$\omega _\theta ^{\prime \prime }=\frac{\mathrm{sin}\theta }{\left(1+\mathrm{cos}^2\theta \right)^{3/2}}=i(1\alpha ^2)\sqrt{2\alpha ^21}.$$ (54) Since we know that $`\mathrm{cos}^2\theta _\alpha =\alpha ^2/(1\alpha ^2),`$ we will obtain $$1+\mathrm{cos}^2\theta _\alpha =1/(1\alpha ^2),$$ (55) and we already have (see equation (44)) $$e^{i\theta _\alpha }=\frac{\sqrt{1\alpha ^2}}{\alpha +\sqrt{2\alpha ^21}}.$$ (56) The standard formula from steepest descents tells us that $$\psi _R(n,t)\frac{(1)^{n+1}}{2\pi }\frac{e^{i\theta _\alpha i(\omega _\alpha \theta _\alpha \alpha )t}}{\sqrt{1+\mathrm{cos}^2\theta _\alpha }}\sqrt{\frac{2\pi }{t\|\omega _\alpha ^{^{\prime \prime }}\|}}.$$ (57) The theorem then follows by using equations (50), (53), (55) and (56). ∎ Remark: The wave-mechanics calculation is conceptually much simpler than the path-integral analysis in Carteret et al. . It is also simpler than the calculations of Chen and Ismail . ## 3 Equivalence of the two approaches We have now completed the calculation begun in and obtained uniformly convergent asymptotics for the wavefunction via both methods. However, the functions $`\psi _L`$ and $`\psi _R`$ derived by each route did not appear to be the same. If they really were different, this would be very alarming, as it would imply either that there is something wrong with Feynman’s equivalence argument in or, more likely, that there was something wrong with our calculation! The first thing to note is that the quantity raised to the power $`t`$ in equation (52) dominates the asymptotics of the logarithm of the functions $`\psi _\mathrm{L}`$ and $`\psi _\mathrm{R}`$ from the Schrödinger representation. Let us call it $`\stackrel{~}{B}(\alpha );`$ that is, $$\stackrel{~}{B}(\alpha )=\left(\frac{\sqrt{1\alpha ^2}}{\alpha +\sqrt{2\alpha ^21}}\right)^\alpha \frac{1+\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}.$$ (58) These estimates agree with the asymptotics obtained using the method of Saff and Varga as used in , although this is not yet apparent. According to the calculation in , the corresponding quantity from the path-integral representation, namely $`B(\alpha ),`$ should be $$2^{\frac{\alpha }{2}}\times \left(\frac{1+2\alpha \sqrt{2\alpha ^21}}{1+\alpha }\right)^\alpha \left(\frac{\alpha ^2+\sqrt{2\alpha ^21}}{1\alpha ^2}\right)^{(1\alpha )/2},$$ (59) which would seem to be a different function. The demonstration of the equivalence to the result obtained by Saff and Varga’s method needs some identities, starting with $$\left(1+\sqrt{2\alpha ^21}\right)^2=2\left(\alpha ^2+\sqrt{2\alpha ^21}\right).$$ (60) Combining $`1+2\alpha \sqrt{2\alpha ^21}`$ with the other quantities is rather fiddly (see below). To show that the two solutions (58) and (59) are equivalent, we will now employ the identity $$\frac{1+2\alpha \sqrt{2\alpha ^21}}{1+\alpha }=\frac{1+\sqrt{2\alpha ^21}}{\alpha +\sqrt{2\alpha ^21}},$$ (61) which can easily be verified by cross-multiplication. It follows immediately from this identity that $$\left(\frac{1+2\alpha \sqrt{2\alpha 1}}{1+\alpha }\right)^\alpha =\left(\frac{1+\sqrt{2\alpha ^21}}{\alpha +\sqrt{2\alpha ^21}}\right)^\alpha $$ (62) and from equation (60) that $`\left({\displaystyle \frac{1\alpha ^2}{\alpha ^2+\sqrt{2\alpha ^21}}}\right)^{\frac{\alpha ^21}{2}}`$ $`=\left({\displaystyle \frac{2(1\alpha ^2)}{(1+\sqrt{2\alpha ^21})^2}}\right)^{\frac{\alpha 1}{2}}`$ (63) $`=\left({\displaystyle \frac{\sqrt{2}\sqrt{1\alpha ^2}}{1+\sqrt{2\alpha ^21}}}\right)^\alpha {\displaystyle \frac{1+\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}}.`$ (64) Thus $$\begin{array}{c}\left(\frac{1+2\alpha \sqrt{2\alpha ^21}}{1+\alpha }\right)^\alpha \left(\frac{1\alpha ^2}{\alpha ^2\sqrt{2\alpha ^21}}\right)^{\frac{\alpha 1}{2}}\hfill \\ \hfill =2^{\alpha /2}\left(\frac{\sqrt{1\alpha ^2}}{\alpha +\sqrt{2\alpha ^21}}\right)^\alpha \frac{1+\sqrt{2\alpha ^21}}{\sqrt{2}\sqrt{1\alpha ^2}}\end{array}$$ (65) so that the formulæ for $`\stackrel{~}{B}(\alpha )`$ and $`B(\alpha ),`$ in (58) and (59) respectively are indeed equivalent. In fact, the representations for the two functions are completely equivalent, but some of the steps required to prove this need some rather subtle calculations involving special functions, as we will now demonstrate. These technical lemmas will also allow us to rederive the symmetry relations in the wave-mechanics picture that were first proved (for the path-integral picture) in . For convenience we recall the integral representations of Theorem 3 $$\psi _L(n,t)=\frac{1}{2\pi }_\pi ^\pi \left(1+\frac{\mathrm{cos}\theta }{\sqrt{1+\mathrm{cos}^2\theta }}\right)e^{i(\omega _\theta +\theta \alpha )t}𝑑\theta ,$$ (66) $$\psi _R(n,t)=\frac{1}{2\pi }_\pi ^\pi \frac{e^{i\theta }}{\sqrt{1+\mathrm{cos}^2\theta }}e^{i(\omega _\theta +\theta \alpha )t}𝑑\theta .$$ (67) where $`\alpha =n/t`$. Note that we have omitted the factors $`\frac{1+(1)^{n+t}}{2}`$ since $`n`$ and $`t`$ must have the same parity because the walker must move at each time-step. From these integral representations obtained in the Schrödinger picture , we prove in this section the symmetry relations for $`\psi _L`$ and $`\psi _R`$ and the relations first proved for the Jacobi polynomials. These results could previously only be proved in the path-integral picture . ### 3.1 Symmetry properties ###### Lemma 2. The function $`\psi _R`$ of (67) satisfies the symmetry relation $$\psi _R(n,t)=(1)^{n+1}\psi _R(n+2,t).$$ (68) This is the symmetry relation in Theorem 2 for $`\psi _R(n,t).`$ ###### Proof. We have, from equation (67) $$\psi _R(n,t)=\frac{1}{\pi }_0^\pi \frac{\mathrm{cos}((n1)\theta )\mathrm{cos}(\omega _\theta t)}{\sqrt{1+\mathrm{cos}^2\theta }}𝑑\theta \frac{1}{\pi }_0^\pi \frac{\mathrm{sin}((n1)\theta )\mathrm{sin}(\omega _\theta t)}{\sqrt{1+\mathrm{cos}^2\theta }}𝑑\theta .$$ (69) The first integral vanishes when $`n`$ is even, the second one when $`n`$ is odd. To verify this, split $`[0,\pi ]=[0,\frac{1}{2}\pi ][\frac{1}{2}\pi ,\pi ]`$ and write on the second interval $`\theta =\pi \theta ^{}`$. We conclude that $`\psi _R(1n,t)=\psi _R(1+n,t)`$ when $`n`$ is odd, and $`\psi _R(1n,t)=\psi _R(1+n,t)`$ when $`n`$ is even. This proves the lemma. ∎ ###### Lemma 3. The function $`\psi _L`$ of equation (66) satisfies the symmetry relation $$(tn)\psi _L(n,t)=(1)^n(t+n)\psi _L(n,t).$$ (70) This is the symmetry relation in Theorem 2 for $`\psi _L(n,t).`$ ###### Proof. We have from (66) $$\psi _L(n,t)=\stackrel{~}{\psi }_L(n,t)+\frac{1}{2\pi }_\pi ^\pi \frac{\mathrm{cos}\theta }{\sqrt{1+\mathrm{cos}^2\theta }}e^{i(t\omega _\theta +n\theta )}𝑑\theta ,$$ (71) where $`\stackrel{~}{\psi }_L(n,t)`$ $`={\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }e^{i(t\omega _\theta +n\theta )}𝑑\theta `$ (72) $`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\mathrm{cos}(n\theta )\mathrm{cos}(t\omega _\theta )𝑑\theta {\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\mathrm{sin}(n\theta )\mathrm{sin}(t\omega _\theta )𝑑\theta .`$ (73) The first integral vanishes when $`n`$ is odd, the second one when $`n`$ is even. This gives the symmetry relation $$\stackrel{~}{\psi }_L(n,t)=(1)^n\stackrel{~}{\psi }_L(n,t).$$ (74) Next observe that $$\frac{d\omega _\theta }{d\theta }=\frac{\mathrm{cos}\theta }{\sqrt{1+\mathrm{cos}^2\theta }},$$ (75) and that an integration by parts in the integral in equation (71) gives us that $$\psi _L(n,t)=\frac{tn}{t}\stackrel{~}{\psi }_L(n,t).$$ (76) Finally we use equations (74) and (76) to complete the proof of the lemma. ∎ Remark: The symmetry relations for $`\psi _L`$ and $`\psi _R`$ also follow from the following property of the Jacobi polynomials: $$\left(\genfrac{}{}{0pt}{}{m}{\mathrm{}}\right)J_m^{(u,\mathrm{})}(x)=\left(\genfrac{}{}{0pt}{}{m+u}{\mathrm{}}\right)\left(\frac{1+x}{2}\right)^{\mathrm{}}J_m\mathrm{}^{(u,\mathrm{})}(x),0\mathrm{}m.$$ (77) This formula follows from the representation of the Jacobi polynomial in terms of the hypergeometric function (cf. \[34, p. 151, (6.35)\]) in combination with a functional relation of this function (third line of equation (5.5) in \[34, p. 110\]). The result in (77) combines the first case in (10) with the second case, and it also implies the symmetry rule for $`\psi _R,`$ and similarly for (11). ### 3.2 The $`𝝍`$functions in terms of Jacobi polynomials We will now prove that the $`\psi `$functions with the integral representations given in (66) and (67) can be written in terms of the Jacobi polynomials as in Theorem 2. This step will require the use of generating functions. These provide a method for writing a series as the coefficients of a formal power series, in a dummy variable, $`z,`$ for ease of manipulation. The powers of $`z`$ are then the summation variable for the series. For a basic introduction to the theory of generating functions, see ; for an advanced treatment, see . We will use these generating functions to give us an exact representation of the $`\mathrm{\Psi }`$ functions as functions of $`t,`$ thus each labelled term in the series will be the function for that value of $`t.`$ Therefore, our approach will be based on generating functions that contain the $`\psi `$functions with $`t`$ as the summation variable. We only consider sums of $`\psi `$functions with $`n`$ and $`t`$ having the same parity. ### 3.3 Some generating functions for $`𝝍`$ For convenience, we will define the following generating functions, which can be obtained from the Schrödinger representation of the wavefunction. ###### Theorem 5. Consider the generating functions for $`\|z\|<1`$: $`F_m(z)`$ $`={\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}\psi _R(2m+1,2t+1)z^t,`$ (78) $`G_m(z)`$ $`={\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}\psi _R(2m,2t)z^t,`$ (79) $`H_m(z)`$ $`={\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}\stackrel{~}{\psi }_L(2m+1,2t+1)z^t,`$ (80) $`I_m(z)`$ $`={\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}\stackrel{~}{\psi }_L(2m,2t)z^t,`$ (81) where $`\stackrel{~}{\psi }_L(n,t)`$ is defined in (72). After the summations have been performed these functions become, respectively, $`F_m(z)`$ $`={\displaystyle \frac{2^{m\frac{1}{2}}z^m}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})^{2m}}},`$ $`m=0,1,2,\mathrm{},`$ (82) $`G_m(z)`$ $`={\displaystyle \frac{2^{m1}z^m}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})^{2m1}}},`$ $`m=1,2,3\mathrm{},`$ (83) $`G_0(z)`$ $`={\displaystyle \frac{z}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})}},`$ (84) $`H_m(z)`$ $`={\displaystyle \frac{2^{m\frac{1}{2}}(1+z)z^m}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})^{2m+1}}},`$ $`m=0,1,2,\mathrm{}`$ (85) $`I_m(z)`$ $`={\displaystyle \frac{2^{m1}(1+z)z^m}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})^{2m}}},`$ $`m=1,2,3\mathrm{},`$ (86) $`I_0(z)`$ $`={\displaystyle \frac{1}{\sqrt{1+z^2}}}{\displaystyle \frac{z}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})}}.`$ (87) These summations can be done using some fiddly, but essentially mechanical manipulations; we have included detailed proofs for $`F_m(z)`$ and a sketch of that for $`G_m(z)`$ in an appendix, in Subsection 6.1. The other generating functions can be obtained via similar constructions. For $`H_m(z)`$ we obtain $`H_m(z)`$ $`={\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }e^{i((2m+1)\theta +\omega _\theta )}{\displaystyle \frac{1}{1ze^{2i\omega _\theta }}}𝑑\theta `$ (88) $`={\displaystyle \frac{1+z}{2\pi \sqrt{2}}}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}((2m+2)\theta )\mathrm{cos}(2m\theta )}{12z\mathrm{cos}2\omega _\theta +z^2}}𝑑\theta `$ (89) $`={\displaystyle \frac{1+z}{2(1z)}}\left[F_{m+1}(z)F_m(z)\right].`$ (90) Taken together with (82), this becomes equation (85). For $`I_m(z)`$ we have a useful intermediate result, namely that $$I_m(z)=\frac{\sqrt{2}}{4(1z)}\left[2(2z)F_m(z)zF_{\|m1\|}(z)zF_{m+1}(z)\right].$$ (91) It remains to compare these with some generating functions for Jacobi polynomials. For an introduction to the theory of Jacobi polynomials we refer the reader to and . This correspondence between the two sets of generating functions forms the central plank of the Feynman equivalence between the two representations of the wave-function for this system. ### 3.4 Comparing the generating functions for $`𝝍`$ We now compare the generating functions (82) – (87) from the wave-mechanics representation, with the generating function of the Jacobi polynomials (cf. \[7, p. 298\]) from the path-integral representation to complete our proof of the equivalence of the two approaches: $$\underset{k=0}{\overset{\mathrm{}}{}}J_k^{(r,s)}(x)z^k=\frac{2^{r+s}}{R(1z+R)^r(1+z+R)^s},\|z\|<1,$$ (92) where $`R=\sqrt{12xz+z^2}`$, which for $`x=0`$ becomes $$\underset{k=0}{\overset{\mathrm{}}{}}J_k^{(r,s)}(0)z^k=\frac{2^{r+s}}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})^r(1+z+\sqrt{1+z^2})^s}.$$ (93) By applying the Cauchy integral formula we find that $$J_k^{(r,s)}(0)=\frac{2^{r+s}}{2\pi i}\frac{dz}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})^r(1+z+\sqrt{1+z^2})^sz^{k+1}},$$ (94) where the integral is taken over a circle with radius less than unity. ###### Lemma 4. For $`n=0`$ we have $`\psi _R(0,0)`$ $`=0,`$ (95) $`\psi _R(0,2t)`$ $`={\displaystyle \frac{1}{2}}J_{t1}^{(1,0)}(0)={\displaystyle \frac{1}{2}}(1)^{t1}J_{t1}^{(0,1)}(0),t=2,4,6,\mathrm{}.`$ (96) For $`0<nt`$, $`n`$ and $`t`$ having the same parity, we have $$\psi _R(n,t)=2^{\frac{n}{2}}(1)^{\frac{tn}{2}}(1)^{n+1}J_{\frac{tn}{2}}^{(0,n1)}(0).$$ (97) This gives the first case of equation (10) from the wave-mechanics representation, which was originally obtained via the path-integral method. ###### Proof. From (78) and (82) we find, as in (94), $$\psi _R(2m+1,2t+1)=\frac{2^{m\frac{1}{2}}}{2\pi i}\frac{z^m}{\sqrt{1+z^2}(1z+\sqrt{1+z^2})^{2m}}\frac{dz}{z^{t+1}}.$$ (98) We now compare (94) with (98) and take $`u=2m,`$ $`v=0,`$ and $`s=tm.`$ This gives $$\psi _R(2m+1,2t+1)=2^{m\frac{1}{2}}J_{tm}^{(2m,0)}(0),0mt.$$ (99) Using the symmetry rule for the Jacobi polynomials $$J_n^{(r,s)}(x)=(1)^nJ_n^{(s,r)}(x),$$ (100) which follows from (92) by putting $`xx,zz`$, we find $$\psi _R(2m+1,2t+1)=2^{m\frac{1}{2}}(1)^{tm}J_{tm}^{(0,2m)}(0),0mt.$$ (101) For the even case, we obtain from (79), (83), and (100) $$\psi _R(2m,2t)=2^mJ_{tm}^{(2m1,0)}(0)=2^m(1)^{tm}J_{tm}^{(0,2m1)}(0),0<mt,$$ (102) and from (79) and (84) we obtain (95). Combining (101) and (102) in one formula gives (97). This proves the lemma. ∎ ###### Lemma 5. For $`n=0`$ we have $`\psi _L(0,0)`$ $`=1,`$ (103) $`\psi _L(0,t)`$ $`=J_{t/2}^{(0,0)}(0){\displaystyle \frac{1}{2}}J_{t/21}^{(1,0)}(0),t=2,4,6,\mathrm{}.`$ (104) For $`0<n<t`$, $`n`$ and $`t`$ having the same parity, we have $$\psi _L(n,t)=2^{n/21}(1)^{(tn)/21}J_{(tn)/21}^{(1,n)}(0).$$ (105) This gives the first case of (11) from the path-integral representation, now obtained via wave-mechanics. ###### Proof. We obtain from (80) and (85) $`\stackrel{~}{\psi }_L(2m+1,2t+1)`$ $`=2^{m\frac{3}{2}}\left[J_{tm}^{(2m+1,0)}(0)+J_{tm1}^{(2m+1,0)}(0)\right]`$ (106) $`=2^{m\frac{3}{2}}{\displaystyle \frac{2t+1}{t+m+1}}J_{tm}^{(2m+1,1)}(0),`$ (107) where we have used the relation for the Jacobi polynomials (cf. \[1, p.782, (22.7.19)\]) $$(u+v+2k)J_k^{(u,v1)}(x)=(u+v+k)J_k^{(u,v)}(x)+(u+k)J_{k1}^{(u,v)}(x).$$ (108) We use (77), (76) and (100), and obtain $$\psi _L(2m+1,2t+1)=2^{m\frac{3}{2}}(1)^{tm1}J_{tm1}^{(1,2m+1)}(0).$$ (109) For the even case we obtain from (81) and (86) $`\stackrel{~}{\psi }_L(2m,2t)`$ $`=2^{m1}\left[J_{tm}^{(2m,0)}(0)+J_{tm1}^{(2m,0)}(0)\right]`$ (110) $`=2^{m1}(1)^{tm1}{\displaystyle \frac{t}{tm}}J_{tm1}^{(1,2m)}(0),`$ (111) where we used (108), (77) and (100). By using (76) we obtain $$\psi _L(2m,2t)=2^{m1}(1)^{tm1}J_{tm1}^{(1,2m)}(0),m=1,2,3,\mathrm{}.$$ (112) From (81) and (87) we obtain (103). Combining (109) and (112) into a single formula gives equation (105). This proves the lemma. ∎ As we have now established that both sets of generating functions match, this completes the proof of the equivalence of the results obtained via the path-integral and wave-mechanics representations. It should be noted that we have proved the two representations are exactly equivalent for all time, as opposed to being only asymptotically equivalent in the long-time limit. It is one of the curious features of generating function methods that they can be used to prove the existence of a one-to-one correspondence between the two sets (i.e., our $`\psi `$-functions) counted by the two series, without actually finding the explicit bijection. ### 3.5 Summary of the equivalence results While the coined quantum walk can be thought of as a quantum analogue of the discrete-time classical random walk , it should be noted that the quantum model inherits its discrete time parameter directly from the classical model; the discreteness was not introduced by hand as part of the quantization procedure. Also, we have not defined a Hamiltonian for this system at all, so the problem of ambiguities in the time derivatives of the action does not arise. We have now shown the full Feynman equivalence for this system, though some results seem easier to derive in one approach than in the other. * We have obtained the symmetry rules directly from the integral representations for the $`\psi `$-functions. It is not necessary to represent the $`\psi `$-functions as Jacobi poynomials and then use the symmetry properties of Jacobi polynomials (as was done in ) to obtain this result. * The relations between the $`\psi `$-functions and the Jacobi polynomials have been obtained directly from the integral representations, though we needed to develop some technical tools for this. These consisted of a few extra properties of the Jacobi polynomials, and the generating functions containing the $`\psi `$-functions. The proofs using these generating functions are conceptually straightforward, although the details are quite technical. The Feynman path-integral approach of (which is a finite sum here) would seem to be simpler if these relations are all one wants. * We were able to establish some new expressions for the values of the two components of the $`\psi `$-function at $`n=0`$ as a function of time, in equations (95) and (103), which were not known before. ## 4 Physical interpretation of these results Now that we have established that the rather counter-intuitive results obtained from the wave-mechanics picture really are equivalent to the Airy functions obtained from the path-integral approach, we are left with the little mystery of their physical interpretation. In this region of exponential decay these waves have complex wavenumbers. This phenomenon is called evanescence. And herein lies the mystery; the conventional wisdom is that evanescent waves are only ever seen in the presence of absorbing media, such as light waves being absorbed into a conducting surface, but there is no such surface here and the evolution is unitary, by the initial assumptions that went into constructing the model. In fact, the phenomenon of evanescence is rather more widespread; it occurs in a great many systems if you know where to look. In a pioneering paper in the early 1990s, Michael Berry showed that evanescent behaviour is much more common than had been previously thought, after being inspired by some work by Aharonov, Anandan, Popescu and Vaidman in . Berry gave a detailed discussion of how this phenomenon occurs in optics, at the edges of the almost ubiquitous “Gaussian” beams in . The wavefunction for this system tends to an Airy function in the asymptotic limit, as was proved analytically in . We evaluated the integrals in the path-integral picture using the method of steepest descents, which in this case featured a pair of coalescing saddle-points . Since then, various authors have discussed the connection between the discrete walk on the infinite line and interference phenomena in the quantum optics of dispersive media. This connection was first described by Knight, Roldan and Sipe in a series of papers and further clarified by Kendon and Sanders in . As we enter the exponential decay region, the two original stationary points of the phase function merge and then two new saddle-points are born, which move off the real axis as a complex conjugate pair. The behaviour of the momentum closely follows that of $`\omega _\theta ,`$ which was plotted in Figure 2. In the wave-mechanics picture, we found that the momentum becomes purely imaginary in the exponential decay region; indeed, the techniques we used to evaluate the integral relied on this fact. So, the behaviour of the walk in the exponential decay region is a pure exponential decay; there is no oscillatory behaviour. Within the interpretation begun by Knight et al., it was first suggested to us by Achim Kempf that the specific evanescence phenomena that we have discussed in this paper are analogues of what are called the Sommerfeld and Brillouin precursors. Specifically, the exponential decay region can be identified with the Sommerfeld precursor (see for example, ) and the distinctive peaks in the probability distribution would be an example of the Brillouin Precursor (see for example ). Our results here and the previous results in provide the first analytic evidence for this identification. ## 5 Conclusions In this paper we have completed the analysis begun in , thus meeting the challenge made in to prove all their theorems about the unrestricted quantum walk on the line in both the path-integral and wave-mechanics representations. We have also proved some additional identities that we believe to be novel. In the course of doing this, we have had to generalise the method of stationary phase in a way that may have applications beyond this problem. We have also proved the exact Feynman equivalence between the two representations directly, by reducing the problem to purely combinatorical constructions which may also be of wider interest. Lastly, we have supplied a physical interpretation for our results, in terms of certain evanescent phenomena from the quantum optics of dispersive media. This interpretation is a somewhat counter-intuitive one, as it would seem to require an effective dissipation that acts on the walker in a way that is analogous to the effect of a dielectric medium on light, despite the fact that the evolution of the system is unitary by assumption. ### Acknowledgments HAC would like to acknowledge some inspiring conversations with Sir Michael Berry, Mourad Ismail and Achim Kempf. HAC was supported by MITACS, and would like to thank the Perimeter Institute and the IQC at the University of Waterloo for hospitality. LBR would also like to thank Ashwin Nayak for some interesting conversations. The research of LBR was partially supported by an NSERC operating grant. NMT acknowledges financial support from Ministerio de Educación y Ciencia (Programa de Sabáticos) from project SAB2003-0113. ## 6 Appendices ### 6.1 Construction of the generating functions Here we give the details of the construction for the generating functions from Subsection 3.3. #### 6.1.1 Proof of the construction for $`𝑭_𝒎\mathbf{(}𝒛\mathbf{)}`$ ###### Proof. We give a detailed proof for $`F_m(z)`$. We substitute equation (67) into equation (78) and obtain $$\begin{array}{c}F_m(z)=\frac{1}{2\pi }_\pi ^\pi \frac{e^{i(2m\theta +\omega _\theta )}}{\sqrt{1+\mathrm{cos}^2\theta }}\frac{1}{1ze^{2i\omega _\theta }}𝑑\theta \hfill \\ \hfill =\frac{1}{2\pi }_\pi ^\pi \frac{e^{i(2m\theta +\omega _\theta )}}{\sqrt{1+\mathrm{cos}^2\theta }}\frac{1z\mathrm{cos}2\omega _\theta iz\mathrm{sin}2\omega _\theta }{12z\mathrm{cos}2\omega _\theta +z^2}𝑑\theta \\ \hfill =\frac{1}{2\pi }_\pi ^\pi \frac{\{\mathrm{cos}(2m\theta +\omega _\theta )i\mathrm{sin}(2m\theta +\omega _\theta )\}\{1z\mathrm{cos}2\omega _\theta iz\mathrm{sin}2\omega _\theta \}}{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}𝑑\theta .\end{array}$$ (113) Since $`\omega _\theta `$ is an odd function, the imaginary parts will vanish. The interval $`[\pi ,\pi ]`$ can be reduced to $`[0,\pi ]`$ and we obtain $`F_m(z)`$ $`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}(2m\theta +\omega _\theta )(1z\mathrm{cos}2\omega _\theta )z\mathrm{sin}2\omega _\theta \mathrm{sin}(2m\theta +\omega _\theta )}{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta `$ (114) $`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}(2m\theta +\omega _\theta )}{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta `$ $`{\displaystyle \frac{z}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}(2m\theta +\omega _\theta )\mathrm{cos}(2\omega _\theta )+\mathrm{sin}(2m\theta +\omega _\theta )\mathrm{sin}2\omega _\theta }{\sqrt{1+\mathrm{cos}^2(\theta )}(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta .`$ (115) The numerator of the second integral can be written as $`\mathrm{cos}(2m\theta \omega _\theta )`$, and it follows that $`F_m(z)`$ $`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}(2m\theta +\omega _\theta )}{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta `$ $`{\displaystyle \frac{z}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}(2m\theta \omega _\theta )}{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta .`$ (116) We can now use simple trigonometric identities for the cosines in the numerators, to obtain $`F_m(z)`$ $`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}2m\theta \mathrm{cos}\omega _\theta \mathrm{sin}2m\theta \mathrm{sin}\omega _\theta }{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta `$ $`{\displaystyle \frac{z}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}2m\theta \mathrm{cos}\omega _\theta +\mathrm{sin}2m\theta \mathrm{sin}\omega _\theta }{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta .`$ (117) Now we use $`\mathrm{sin}\omega _\theta =(\mathrm{sin}\theta )/\sqrt{2}`$, and observe that the terms with the sine functions do not contribute to the integrals; this follows easily by performing the transformation $`\theta =\theta ^{}+\frac{1}{2}\pi `$. Using also $`\mathrm{cos}\omega _\theta =\sqrt{\frac{1}{2}(1+\mathrm{cos}^2\theta )}`$, we obtain $$F_m(z)=\frac{1z}{\pi \sqrt{2}}_0^\pi \frac{\mathrm{cos}(2m\theta )}{12z\mathrm{cos}2\omega _\theta +z^2}𝑑\theta .$$ (118) For the final step we use formula (3.615) (1) of Gradshteyn and Ryzhik , that is, $$_0^{\frac{1}{2}\pi }\frac{\mathrm{cos}(2m\theta )}{1a^2\mathrm{sin}^2\theta }𝑑\theta =\frac{(1)^m\pi }{2\sqrt{1a^2}}\frac{(1\sqrt{1a^2})^{2m}}{a^{2m}},\|a^2\|<1,m=0,1,2,\mathrm{}.$$ (119) We observe that $`\mathrm{cos}2\omega _\theta =\mathrm{cos}^2\theta `$, and take $`a^2=2z/(1z)^2`$. This gives the expression in (82), as advertised. #### 6.1.2 Outline of the proof for $`𝑮_𝒎\mathbf{(}𝒛\mathbf{)}`$ The proof of equation (83) for $`G_m(z)`$ uses essentially the same manipulations: $`G_m(z)`$ $`={\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }{\displaystyle \frac{e^{i(12m)\theta }}{\sqrt{1+\mathrm{cos}^2\theta }}}{\displaystyle \frac{1}{1ze^{2i\omega _\theta }}}𝑑\theta `$ (120) $`={\displaystyle \frac{1}{2\pi }}{\displaystyle _\pi ^\pi }{\displaystyle \frac{\mathrm{cos}(2m1)\theta (1z\mathrm{cos}2\omega _\theta )z\mathrm{sin}(2m1)\theta \mathrm{sin}2\omega _\theta }{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta .`$ (121) The second integral can be broken into two terms as follows: $`G_m(z)`$ $`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}(2m1)\theta }{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta `$ (122) $`{\displaystyle \frac{z}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{cos}(2m1)\theta \mathrm{cos}2\omega _\theta +\mathrm{sin}(2m1)\theta \mathrm{sin}2\omega _\theta }{\sqrt{1+\mathrm{cos}^2\theta }(12z\mathrm{cos}2\omega _\theta +z^2)}}𝑑\theta .`$ (123) The first integral vanishes (substitute $`\theta =\theta ^{}+\frac{1}{2}\pi `$). The same holds for the contributions from the cosine terms in the second integral. This gives $$G_m(z)=\frac{z}{\pi }_0^\pi \frac{\mathrm{sin}(2m1)\theta \mathrm{sin}\theta }{12z\mathrm{cos}2\omega _\theta +z^2}𝑑\theta .$$ (124) Using the fact that $`\mathrm{cos}2m\theta `$ $`=\mathrm{cos}(2m1)\theta \mathrm{cos}\theta +\mathrm{sin}(2m1)\theta \mathrm{sin}\theta `$ (125) $`\mathrm{cos}(2m2)\theta `$ $`=\mathrm{cos}(2m1)\theta \mathrm{cos}\theta \mathrm{sin}(2m1)\theta \mathrm{sin}\theta `$ (126) we can then write $$G_m(z)=\frac{z}{2\pi }_0^\pi \frac{\mathrm{cos}2m\theta \mathrm{cos}(2m2)\theta }{12z\mathrm{cos}2\omega _\theta +z^2}𝑑\theta ,$$ (127) and we see (cf. (118)) that we can express $`G_m(z)`$ in terms of two $`F_m(z)`$ functions. This gives equation (83). ∎ ### 6.2 Lagrange inversion asymptotics Suppose we have an unknown function $`w`$ that we assume can be written as an (as yet unknown) power series. All we know about this power series is that it can be written as a recursion relation $`w=z\phi (w),`$ $`\phi (0)0,`$ (128) where $`\phi (w)`$ is some generating function that is defined as an implicit function of $`z`$. We would like to find $`w`$ as an explicit function of $`z,`$ so we can express some other function $`f(w)`$ as an explicit power series in $`z.`$ Lagrange Inversion enables us to do this, and tells us that we can write $$f(w)=\underset{n1}{}\frac{t^n}{n}[\lambda ^{n1}]f^{^{}}(\lambda )\varphi ^n(\lambda ),$$ (129) where $`\lambda `$ is a dummy variable, denotes differentiation with respect to $`\lambda `$ and the square brackets $`[x^n]`$ is the “Goulden-Jackson” notation for the coefficient of the term in $`x^n.`$ Here are two very simple examples to illustrate the use of Lagrange Inversion. Suppose we were given $`w`$ defined only to be a solution of the equation $$w=ze^w,$$ (130) and we know nothing else about it. Suppose we just want to obtain a power series for $`f(w)=w,`$ where $`\phi (w)=e^w.`$ Then the formula above becomes $$w=\underset{n1}{}\frac{t^n}{n}[\lambda ^{n1}]e^{n\lambda }=\underset{n1}{}\frac{n^{n1}}{n!}.$$ (131) It follows from Stirling’s formula for the factorial that this series converges for $`\|z\|<1/e`$. A slightly less simple example occurs if we are interested in the function $`g(w)=w^2.`$ Then equation (131) would become $$w^2=\underset{n1}{}\frac{t^n}{n}[\lambda ^{n1}]2\lambda e^{n\lambda }=2\underset{n1}{}\frac{t^n}{n}[\lambda ^{n2}]e^{n\lambda }=2\underset{n2}{}t^n\frac{n^{n3}}{(n2)!}.$$ (132) #### 6.2.1 Lagrange inversion for Jacobi polynomials The formula of Lagrange most useful for Jacobi polynomials is from $$\frac{f(\lambda )}{1z\phi ^{}(\lambda )}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{z^n}{n!}\frac{d^n}{dx^n}\{f(x)[\phi (x)]^n\}.$$ (133) where $`\phi `$ was defined in equation (128) and $`x`$ is another dummy variable. The asymptotics of the Jacobi polynomials have been discussed previously by Chen and Ismail and Ambainis et al. used those results to derive some results on the asymptotics of the $`\psi `$-functions. It should be noted that the Chen-Ismail results used the method of Darboux and are not uniform over the full range of $`\alpha `$. For the rest of this section we will briefly discuss Lagrange inversion and show how to derive the integral representations that Carteret et al. used for the $`\psi `$-function to derive uniformly convergent asymptotics. Chen and Ismail’s work on the Jacobi polynomials uses the generating function for $`J_j^{(\gamma +aj,\beta +bj)}(0)`$ of Srivastava and Singhal which in our notation becomes $$\underset{j=0}{\overset{\mathrm{}}{}}J_j^{(\gamma +aj,\beta +bj)}(0)z^j=(1+𝐮)^{\gamma +1}(1+𝐯)^{\beta +1}[1a𝐮b𝐯(1+a+b)\mathrm{𝐮𝐯}]^1,$$ (134) where $`𝐮=𝐯`$ and $`𝐮`$ is implicitly defined as a function of $`z`$ by $$𝐮=\frac{z}{2}(1𝐮)^{1+a}(1+𝐮)^{1+b},$$ (135) In order to obtain the asymptotics, we will need to interpret equation (134) in terms of equation (133), again following Srivastava and Singhal . In the case of $`\psi _L,`$ we should let (see ) $`a=\gamma =0,`$ $`\beta ={\displaystyle \frac{1+\alpha }{1\alpha }},`$ $`b={\displaystyle \frac{2\alpha }{1\alpha }}.`$ (136) Then we should define $$\phi (\lambda )=\frac{\lambda ^21}{2}(1+\lambda )^{2\alpha /(1\alpha )}$$ (137) and $`f`$ should be defined by equation (133), which we must set equal to $$f(\lambda )=(1+\lambda )^{(1+\alpha )/(1\alpha )},$$ (138) following the method originally developed in . We omit the details, since we only wish to use the result. Then $$J_{(tn)/21}^{(0,n+1)}(0)=J_m^{(0+0m,\frac{1+\alpha }{1\alpha }+\frac{2\alpha m}{1\alpha })}(0),$$ (139) and so $$m=(1\alpha )t/21,$$ (140) as required. We would like to obtain an integral representation of the coefficients in equation (134). If $`f(\lambda )`$ and $`\phi (\lambda )`$ are analytic then this can be done using the Cauchy integral formula, thus $$[\lambda ^n]f(\lambda )\varphi ^n(\lambda )=\frac{1}{2\pi i}_Cf(\lambda )\varphi ^n(\lambda )\lambda ^{n1}𝑑\lambda $$ (141) where $`C`$ is a sufficiently small contour around the origin. Equation (141) is an example of a Rodrigues formula for a set of orthogonal polynomials. Another example of these appears in the method for generating orthogonal polynomials using Gram-Schmidt orthogonalization . So we obtain the integral representation for the Jacobi polynomials $$J_n^{(0,2\alpha n/(1\alpha )+\beta )}(0)=\frac{1}{2\pi i}_C\frac{(1+\lambda )^{(1+\alpha )/(1\alpha )}}{\lambda }\left(\frac{\lambda ^21}{2\lambda }(1+\lambda )^{2\alpha /(1\alpha )}\right)^n𝑑\lambda ,$$ (142) as used in . This is the contour integral that Saff and Varga estimated using steepest descents. This is discussed in complete detail in so there is no need to say more here about this particular example. This example shows however that expressing a coefficient obtained using Lagrange inversion as a contour integral in this way and using steepest descents may lead to uniform asymptotic expansions over a wide domain. See the book by Andrews et al. for examples of Lagrange inversion arising in special functions.
warning/0506/hep-ph0506058.html	ar5iv	text	# Axial and tensor charge of the nucleon in the Dirac orbital model ## Abstract Using the expansion of the baryon wave function in a series of products of single quark bispinors (Dirac orbitals), the nonsinglet axial and tensor charge of the nucleon are calculated. The leading term yields $`G_A/G_V=1.27`$ and in good agreement with experiment. Calculation is essentially parameter-free and depends on the string tension $`\sigma `$ and $`\alpha _s`$, fixed at standard values. The importance of lower Dirac bispinor component, yielding 18% to the wave function normalization is stressed. Axial and tensor charges of nucleons are important to characterize the basic structure of the nucleon as composed of strongly coupled quarks -. It is known that nonrelativistic quark models predict $`\frac{G_A}{G_V}=\frac{5}{3}`$ in strong disagreement with experimental value 1.27, while for massless relativistic quarks, e.g. in the MIT bag model, one obtains much smaller values $`\frac{G_A}{G_V}=1.09`$. Thus the calculation of $`\frac{G_A}{G_V}`$ (and tensor charge $`\delta q`$) gives a clue to the relativistic dynamics of quarks in the nucleon. Moreover, in a recent paper it was shown that the knowledge of the ratio of $`\frac{G_A}{G_V}`$ for baryon decays is important for the accurate determination of the CKM matrix element $`V_{us}`$ These considerations justify the systematic analysis of baryon decays in the framework of Dirac Orbital Expansion (DOE) the first part of which is reported below. The contribution of vector and axial hadronic currents, $`V_\mu =i\overline{\psi }_u\gamma _\mu \psi _d`$, $`A_\mu =i\overline{u}\gamma _\mu \gamma _5\psi _d`$, to the neutron $`\beta `$-decay is characterized by the ratio $$\frac{G_A}{G_V}=\frac{p_{}\left\|A_z\right\|n_{}}{p_{}\left\|V_0\right\|n_{}},$$ (1) where $`p_\lambda ,n_\lambda ,\lambda =\pm \frac{1}{2}`$ are proton and neutron wave functions with spin projection $`\lambda `$ -. In a similar way the tensor charge is expressed through the proton matrix element of the tensor current $`T_{\mu \nu }\overline{\psi }i\sigma _{\mu \nu }\gamma _5\psi `$ . To construct the baryon wave function, one starts with the Hamiltonian obtained in the instantaneous approximation from the general Bethe–Salpeter equation: $$\widehat{H}\mathrm{\Psi }(𝐫_1,𝐫_2,𝐫_3)=E\mathrm{\Psi },\widehat{H}=\underset{i=1}{\overset{3}{}}\widehat{H}_i+\mathrm{\Delta }H$$ (2) with $$\widehat{H}_i=𝐩_{(i)}𝜶_{(i)}+\beta _{(i)}(m_i+M(𝐫_i))$$ (3) where $`M(𝐫_i)`$ in the limit of vanishing gluon correlation length is $$M=\sigma \|𝐫_i\|e^{i\gamma _5\widehat{\varphi }(𝐫_i)}$$ and $`𝐫_i=𝐱_i𝐱_0`$, $`𝐱_0`$ is the string-junction coordinate and $`\widehat{\varphi }(𝐫_i)`$ is the Nambu–Goldstone octet. Here $`\mathrm{\Delta }H`$ contains perturbative gluon exchanges. We expand the baryon wave function in a series of products of quark eigenfunctions $`\psi _n^{(i)}=\left(\genfrac{}{}{0pt}{}{v^{(i)}}{w^{(i)}}\right)`$, namely $$\mathrm{\Psi }(𝐫_1,𝐫_2,𝐫_3)=\underset{\{n_i\}}{}\underset{i=1}{\overset{3}{}}\psi _{n_i}^{(i)}(𝐫_i)C_{n_1n_2n_3}$$ (4) In what follows we shall consider the leading valence approximation for the nucleon keeping only the first term in (4), $`\mathrm{\Psi }\mathrm{\Psi }_0`$, which contains the ground state $`S`$–wave Dirac orbitals $`\|u_\lambda `$ and $`\|d_\lambda `$ for $`u`$ and $`d`$ quarks with spins up and down. One has $$\|p=\sqrt{\frac{1}{18}}[2(\|uud+\text{perm.})+(\|uud+\text{perm.})]$$ (5) $$\|n=\sqrt{\frac{1}{18}}[2(\|ddu+\text{perm.})+(\|ddu+\text{perm.})]$$ (6) The expressions (5) and (6) have the same form as in the standard $`SU(4)`$ or $`SU(6)`$ model except for the bispinor contents of $`\|u_\lambda `$ and $`\|d_\lambda `$. Insertion of (5), (6) into (1) yields<sup>1</sup><sup>1</sup>1The sign of $`G_A`$ corresponds to $$\frac{G_A}{G_V}=+\frac{5}{3}\chi _{}\left\|𝚺_3\right\|\chi _{},𝚺=\left(\begin{array}{cc}𝝈& 0\\ 0& 𝝈\end{array}\right)$$ (7) where $`\chi _{}`$ is $$\chi _{}(r,\theta ,\varphi )=\frac{1}{r}\left(\genfrac{}{}{0pt}{}{G(r)\mathrm{\Omega }_{\frac{1}{2},0,\frac{1}{2}}(\theta ,\varphi )}{iF(r)\mathrm{\Omega }_{\frac{1}{2},1,\frac{1}{2}}(\theta ,\varphi )}\right);\underset{0}{\overset{\mathrm{}}{}}\left[G^2(r)+F^2(r)\right]𝑑r=1$$ (8) To take into account perturbative gluon exchange we represent $`\mathrm{\Delta }H`$ effectively as one-particle operators, $$\mathrm{\Delta }H=\underset{i=1}{\overset{3}{}}\left(\frac{\zeta }{r_i}\right)$$ and the equations for $`G(r)`$, $`F(r)`$ acquire the form $$\begin{array}{c}G^{}\frac{1}{r}G\left(E+m+\sigma r+\frac{\zeta }{r}\right)F=0,\\ F^{}+\frac{1}{r}F+\left(Em\sigma r+\frac{\zeta }{r}\right)G=0\end{array}$$ (9) Finally $`G_A`$ can be written as $$G_A=+\frac{5}{3}\left(1\frac{4}{3}\eta \right),\eta =\underset{0}{\overset{\mathrm{}}{}}F^2(r)𝑑r$$ (10) We have computed the values of $`\eta `$ and for $`m=0`$ and two different values of $`\zeta `$: $`\zeta =0`$ and $`\zeta =0.3`$. The results are given in Table 1 Note that for $`m=0`$ $`G_A`$ does not depend on the string tension $`\sigma `$ on dimensional grounds. One can see that in the Table 1 that the resulting $`G_A`$ is in the correct ballpark for $`\zeta [0,0.3]`$. The value $`\zeta =0.3`$ corresponds to the reasonable effective value of $`\alpha _s`$ in the $`qq`$ potential, namely from $`\frac{2}{3}\frac{\alpha _s}{r_{ij}}=\frac{\zeta }{r_i}`$, and $`r_{ij}\sqrt{3}r_i`$, one has $$(\alpha _s)_{\text{eff}}=\frac{3\sqrt{3}}{4}\zeta 0.39$$ and this value of $`\zeta `$ was checked in the actual calculation of the nucleon mass . It is rewarding that the resulting $`G_A=1.27`$ is in close agreement with experiment. Concerning the tensor charge $`\delta q`$, it can be easily calculated the definition of $`\delta q`$, and in the same way as in Eq. (7) one arrives at he expression $$\delta q=\chi _{}\left\|\beta 𝚺_3\right\|\chi _{}$$ (11) As a result one has for $`\zeta =0.3`$: $$\delta q=\delta u\delta d=\frac{5}{3}\left(1\frac{2}{3}\eta \right)1.47$$ (12) Note that in the nonrelativistic limit $`\delta q=G_A`$. One can compare these results with the lattice data , where both $`G_A`$ and $`\delta q`$ are close to each other and are in the interval $`1.12G_A,\delta q1.18`$ for $`m_\pi >0.5\text{ GeV}`$. Calculations of $`\delta q`$ in other methods give results ranging from 1.07 to 1.45, see for refs. and discussion. Note that the anomalous dimension of the tensor charge is small and calculations here and in refer to the scale $`\mu ^2=M_N^2`$. There are two possible unaccounted effects which can influence our results. First, the contribution of other terms in (4) – excited Dirac orbitals. The corresponding multichannel calculations done in for magnetic moments, result in decreasing of the modulus of magnetic moments of proton and neutron by some $`1015\%`$ when one accounts for 4 Dirac orbitals for each quarks, and one can expect the same type of corrections for $`G_A`$. Second, the contribution of chiral degrees of freedom, i.e. of the $`\pi `$, $`\eta `$, $`K`$ exchanges. Again, for nucleon magnetic moments these corrections are typically of the order of $`10\%`$ , and we expect this to be an upper limit for $`G_A`$, since magnetic moments are much more sensitive to the contribution of the lowest Dirac components, than $`G_A`$ and $`\delta q`$, where these contributions enter quadratically and not linearly. These corrections are not taken into account above, which is planned for a subsequent work, where also hyperon semileptonic decays are considered . One should note, that relativistic approach to the baryon wave function, based on the light-cone formalism was shown to improve the qualitative agreement of $`G_A`$ with experiment, however quantitatively still far from experiment. Summarizing, we have calculated nonsinglet axial and tensor charges in the simple relativistic model of the nucleon, where the wave function is a product of three Dirac orbitals of quarks. The resulting value of $`G_A`$ is in excellent agreement with experiment for the choice of the only parameter $`(\alpha _s)_{\text{eff}}=0.39`$, yielding the reasonable value of the nucleon mass. Note that quarks in the baryon prove rather relativistic, so the contribution of the lower quark bispinor component to $`G_A`$ is not negligibly small: $`\eta 0.18`$. This work was supported by the federal program of the Russian Ministry of Industry, Science and Technology # 40.052.1.1.1112, by the Grant for support of Leading Scientific Schools # 1774.2003.2 and in part by the RFBR grant # 03-02-17345.
warning/0506/astro-ph0506704.html	ar5iv	text	# Large non-Gaussianity in multiple-field inflation ## I Introduction The assumption that inflationary fluctuations are Gaussian is a good starting point for the study of cosmological perturbations, but it is only true to linear order in perturbation theory. Since gravity is inherently non-linear, and most inflation models have (self-)interacting potentials, non-linearity must be present at some level in all inflation models. Hence the issue is not whether inflation is non-Gaussian, but how large the non-Gaussianity is. With increasingly precise CMB data becoming available in the near future from the WMAP and Planck satellites and other experiments, we might well hope to detect this non-Gaussianity. This would offer us another key observable to help constrain or confirm specific inflation models and the underlying high-energy theories from which they are derived. As a rough order of magnitude estimate, we note that non-Gaussianity will be detectable by Planck if the bispectrum (the Fourier transform of the three-point correlator) is of the order of the square of the power spectrum komatsu . It follows that to compute the predicted amount of non-Gaussianity in specific inflation models we need to go beyond linear-order perturbation theory. In gp2 ; formalism we introduced a new formalism to deal with the non-linearity during inflation. We will not again summarise the other work dealing with this subject, references for which can be found in formalism or a recent review ngreview . Our formalism is distinguished by being based on a system of fully non-linear equations for long wavelengths, while stochastic sources take into account the continuous influx of short-wavelength fluctuations into the long-wavelength system as the inflationary comoving horizon shrinks. The variables used incorporate both scalar metric and matter perturbations self-consistently and they are invariant under changes of time slicing in the long-wavelength limit. The advantages of our method are threefold: (i) it is physically intuitive and relatively simple to use for quantitative analytic and semi-analytic calculations; (ii) it is valid in a very general multiple-field inflation setting, which includes the possibility of a non-trivial field metric; and (iii) it is well-suited for direct numerical implementation. The first point was already demonstrated in sf , where we computed the non-Gaussianity in single-field slow-roll inflation, while the third point is the subject of a forthcoming paper RSvTnum . The present paper is dedicated to the exploration of the second point, as well as the first. In sf we found, confirming what was known in the literature beforehand (see e.g. maldacena ), that non-Gaussianity in the single-field case is too small to be realistically observable, because it is suppressed by slow-roll factors (actually the scalar spectral index $`n1`$). However, there have been long-standing claims in the literature (see e.g. salopek ; uzan ) that specific multiple-field models can, in principle, create significant non-Gaussianity. Indeed, there has been growing recent interest in models which can produce large primordial non-Gaussianity large ng . A feature shared by most of these models though is that this non-Gaussianity involves some mechanism operating after inflation, usually (p)reheating or later domination of a curvaton field. In this paper we investigate for the first time general multiple-field inflation, not just specific models, presenting a method by which to accurately calculate the resulting three-point correlator during inflation. We find that it is possible to get significant primordial non-Gaussianity without invoking some post-inflationary mechanism even for the simplest two-field quadratic potential. The key mechanism for the production of this large non-Gaussianity is the superhorizon influence of isocurvature perturbations on the adiabatic mode. The former feed into the latter when the background follows a curved trajectory in field space. Note that the example studied in section V.4 illustrates that there is no need for the potential to be interacting. Our aim is to push forward towards a tractable non-Gaussian methodology for the new era of precision cosmology which confronts us. This work is organised as follows. In section II we give the equations from formalism that are used as the starting point for the present investigations. In section III we then derive the general solution for the relevant quantities in multiple-field inflation, culminating in a general expression for the three-point correlator of the adiabatic component of the curvature perturbation, without any slow-roll approximations. This integral solution - equation (29) - is a very useful calculational tool because it gives the three-point correlator entirely in terms of background quantities and linear perturbation quantities at horizon crossing. In section IV we make a leading-order slow-roll approximation to work out the various contributions in the general solution more explicitly. Finally in section V we calculate the bispectrum in an analytic treatment of the two-field case with constant slow-roll parameters. We find that the result can be orders of magnitude larger than for single-field inflation. This result is confirmed with a semi-analytic slow-roll calculation of an explicit model with a quadratic potential in section V.4. Our method yields the full momentum dependence, not just an overall magnitude, and we find that there can be a difference of the order of a few between opposite extreme momentum limits. We conclude in section VI. Parts of this paper, in particular section V, are rather technical, so some readers might be interested in referring to mf2 first, which contains a simplified derivation of only the dominant non-Gaussian contributions in multiple-field inflation, along with a summarised discussion. ## II Multiple-field setup Since in this paper we are explicitly working out the general non-linear multiple-field formalism of formalism , we refer the reader to that paper for derivations and more details of the initial equations. Here we just briefly describe the context and give the relevant equations and definitions to be used as starting point for further calculations. We start from a completely general inflation model, with an arbitrary number of scalar fields $`\varphi ^A`$ (where $`A`$ labels the different fields) and a potential $`V(\varphi ^A)`$ with arbitrary interactions. We also allow for the possibility of a non-trivial field manifold with field metric $`G_{AB}`$. We will consider only scalar modes and make the long-wavelength approximation (i.e. consider only wavelengths larger than the Hubble radius $`1/(aH)`$, where second-order spatial gradients can be neglected compared to time derivatives)<sup>1</sup><sup>1</sup>1Formally this corresponds with taking only the leading-order terms in the gradient expansion. We expect higher-order terms to be subdominant on long wavelengths during inflation, but this statement has only been rigorously verified at the linear level. A calculation to higher order in spatial gradients, or, even better, a full proof of convergence of the expansion, would be desirable. See gradient for more details on the validity of the gradient expansion beyond linear theory.. The spacetime metric $`g_{\mu \nu }`$ and matter Lagrangean are given by $$\mathrm{d}s^2=N^2(t,\text{x})\mathrm{d}t^2+a^2(t,\text{x})\mathrm{d}\text{x}^2,_\mathrm{m}=\frac{1}{2}g^{\mu \nu }_\mu \varphi ^AG_{AB}_\nu \varphi ^BV(\varphi ^A),$$ (1) with $`a`$ the local scale factor and $`N`$ the lapse function. The local expansion rate is defined as $`H\dot{a}/(Na)`$, where the dot denotes a derivative with respect to $`t`$. The proper field velocity is $`\mathrm{\Pi }^A\dot{\varphi }^A/N`$, with length $`\mathrm{\Pi }`$. We also define local slow-roll parameters as $$\stackrel{~}{ϵ}(t,\text{x})\frac{\kappa ^2\mathrm{\Pi }^2}{2H^2},\stackrel{~}{\eta }^A(t,\text{x})\frac{3H\mathrm{\Pi }^A+G^{AB}_BV}{H\mathrm{\Pi }},\stackrel{~}{\xi }^A(t,\text{x})\frac{𝒟_B^AV}{H^2}\frac{\mathrm{\Pi }^B}{\mathrm{\Pi }}+3\stackrel{~}{ϵ}\frac{\mathrm{\Pi }^A}{\mathrm{\Pi }}3\stackrel{~}{\eta }^A,$$ (2) where $`\kappa ^28\pi G=8\pi /m_{\mathrm{pl}}^2`$ and $`𝒟_B`$ is a covariant derivative with respect to the field $`\varphi ^B`$. For the first part of the paper we will not make a slow-roll approximation, and consider these definitions as just a short-hand notation. When we do make this approximation, from section IV, $`\stackrel{~}{ϵ}`$ and $`\stackrel{~}{\eta }^A`$ are first order in slow roll, while $`\stackrel{~}{\xi }^A`$ is second order. Finally, we choose the gauge where $$t=\mathrm{ln}(aH)NH=(1\stackrel{~}{ϵ})^1.$$ (3) In this gauge horizon exit of a mode, $`k=aH`$, occurs simultaneously for all spatial points and calculations are simpler. We will make use of a preferred basis in field space, defined as follows. The first basis vector $`e_1^A`$ is the direction of the field velocity. Next, $`e_2^A`$ is defined as the direction of that part of the field acceleration that is perpendicular to $`e_1^A`$. One continues this orthogonalisation process with higher derivatives until a complete basis is found. Explicit expressions can be found in formalism , here we only give $`e_1^A\mathrm{\Pi }^A/\mathrm{\Pi }`$. Now one can take components of vectors in this basis and we define, for example for $`\zeta _i^A`$ (defined below in (6)) and $`\stackrel{~}{\eta }^A`$: $$\zeta _i^me_{mA}\zeta _i^A,\stackrel{~}{\eta }^{}e_1^A\stackrel{~}{\eta }_A,\stackrel{~}{\eta }^{}e_2^A\stackrel{~}{\eta }_A.$$ (4) Note that, unlike for the index $`A`$, there is no difference between upper and lower indices for the $`m`$. By construction there are no other components of $`\stackrel{~}{\eta }^A`$, so that one can write $`\stackrel{~}{\eta }^{}=\|\stackrel{~}{\eta }^A\stackrel{~}{\eta }^{}e_1^A\|`$. We also define $$Z_{mn}\frac{1}{NH}e_{mA}𝒟_te_n^A,$$ (5) where $`𝒟_t`$ is the covariant time derivative containing the connection of the field manifold. $`Z_{mn}`$ is antisymmetric and only non-zero just above and below the diagonal, and first order in slow roll. Its explicit form in terms of slow-roll parameters can be found in vantent ; here we only need that $`Z_{12}=Z_{21}=\stackrel{~}{\eta }^{}`$. As discussed in gp2 ; formalism it is useful to work with the following combination of spatial gradients to describe the fully non-linear inhomogeneities: $$\zeta _i^A(t,\text{x})e_1^A(t,\text{x})_i\mathrm{ln}a(t,\text{x})\frac{\kappa }{\sqrt{2\stackrel{~}{ϵ}(t,\text{x})}}_i\varphi ^A(t,\text{x})\zeta _i^m\delta _{m1}_i\mathrm{ln}a\frac{\kappa }{\sqrt{2\stackrel{~}{ϵ}}}e_{mA}_i\varphi ^A,$$ (6) which is invariant under changes of time slicing, up to second-order spatial gradients gp ; formalism . Note that, when linearised, $`\zeta _i^1`$ (the $`m`$=1 component of $`\zeta _i^m`$) is the spatial gradient of the well-known $`\zeta `$ from the literature, the curvature perturbation. In formalism we derived a fully non-linear equation of motion for $`\zeta _i^m`$ without any slow-roll approximations: $$\{\begin{array}{c}\dot{\zeta }_i^m\theta _i^m=𝒮_i^m\hfill \\ \dot{\theta }_i^m+\left(\frac{(32\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{}3\stackrel{~}{ϵ}^24\stackrel{~}{ϵ}\stackrel{~}{\eta }^{})\delta _{mn}}{(1\stackrel{~}{ϵ})^2}+\frac{2Z_{mn}}{1\stackrel{~}{ϵ}}\right)\theta _i^n+\frac{\mathrm{\Xi }_{mn}}{(1\stackrel{~}{ϵ})^2}\zeta _i^n=𝒥_i^m\hfill \end{array}$$ (7) where $`\theta _i^m`$ is the velocity corresponding with $`\zeta _i^m`$ and $`\mathrm{\Xi }_{mn}`$ $``$ $`{\displaystyle \frac{V_{mn}}{H^2}}{\displaystyle \frac{2\stackrel{~}{ϵ}}{\kappa ^2}}R_{m11n}+(1\stackrel{~}{ϵ})\dot{Z}_{mn}+Z_{mp}Z_{pn}+\left(32\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{}\stackrel{~}{ϵ}^22\stackrel{~}{ϵ}\stackrel{~}{\eta }^{}\right){\displaystyle \frac{Z_{mn}}{1\stackrel{~}{ϵ}}}`$ $`+\left(3\stackrel{~}{ϵ}+3\stackrel{~}{\eta }^{}+2\stackrel{~}{ϵ}^2+4\stackrel{~}{ϵ}\stackrel{~}{\eta }^{}+(\stackrel{~}{\eta }^{})^2+\stackrel{~}{\xi }^{}\right)\delta _{mn}2\stackrel{~}{ϵ}\left((3+\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{})\delta _{m1}\delta _{n1}+\stackrel{~}{\eta }^{}(\delta _{m1}\delta _{n2}+\delta _{m2}\delta _{n1})\right),`$ where $`V_{mn}e_m^A(𝒟_B_AV)e_n^B`$ and $`R_{m11n}e_m^AR_{ABCD}e_1^Be_1^Ce_n^D`$ with $`R_{BCD}^A`$ the curvature tensor of the field manifold. Although for the first part of the paper we will not make a slow-roll approximation, we give here immediately the leading-order slow-roll approximation of (7), which we will be using in the second part, to show that things simplify considerably in that case: $$\{\begin{array}{c}\dot{\zeta }_i^m\theta _i^m=𝒮_i^m\hfill \\ \dot{\theta }_i^m+3\theta _i^m+\left(\frac{V_{mn}}{H^2}+3Z_{mn}+3(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\delta _{mn}6\stackrel{~}{ϵ}\delta _{m1}\delta _{n1}\right)\zeta _i^n=𝒥_i^m\hfill \end{array}$$ (9) The stochastic source terms $`𝒮_i^m`$ and $`𝒥_i^m`$ are given by $`𝒮_i^m`$ $`=`$ $`{\displaystyle \frac{\kappa }{a\sqrt{2\stackrel{~}{ϵ}}}}{\displaystyle \frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\dot{𝒲}(k)Q_{mn}^{\mathrm{lin}}(k)\alpha _n(\text{k})\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}+\mathrm{c}.\mathrm{c}.,`$ $`𝒥_i^m`$ $`=`$ $`{\displaystyle \frac{\kappa }{a\sqrt{2\stackrel{~}{ϵ}}}}{\displaystyle \frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\dot{𝒲}(k)\left[\dot{Q}_{mn}^{\mathrm{lin}}(k)\delta _{np}Q_{mn}^{\mathrm{lin}}(k)\frac{(1+\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\delta _{np}+Z_{np}}{1\stackrel{~}{ϵ}}\right]\alpha _p(\text{k})\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}+\mathrm{c}.\mathrm{c}.,`$ (10) where c.c. denotes the complex conjugate. The perturbation quantity $`Q_{mn}^{\mathrm{lin}}`$ is the solution from linear theory for the multiple-field generalisation of the Sasaki-Mukhanov variable $`Qa\sqrt{2\stackrel{~}{ϵ}}\zeta /\kappa `$. It can be computed exactly numerically, or analytically within the slow-roll approximation vantent . The $`\alpha _m(\text{k})`$ are a set of Gaussian complex random numbers satisfying $$\alpha _m(\text{k})\alpha _n^{}(\text{k}^{})=\delta _{mn}\delta ^3(\text{k}\text{k}^{}),\alpha _m(\text{k})\alpha _n(\text{k}^{})=0.$$ (11) The quantity $`𝒲(k)`$ is the Fourier transform of an appropriate smoothing window function which cuts off modes with wavelengths smaller than the Hubble radius; we choose it to be a Gaussian with smoothing length $`Rc/(aH)=c\mathrm{e}^t`$, where $`c`$3–5: $$𝒲(k)=\mathrm{e}^{k^2R^2/2}\dot{𝒲}(k)=k^2R^2\mathrm{e}^{k^2R^2/2}.$$ (12) Since $`\zeta _i^m`$ and $`\theta _i^m`$ are smoothed long-wavelength variables, the appropriate initial conditions are that they should be zero at early times when all the modes are sub-horizon. Hence, $$\underset{t\mathrm{}}{lim}\zeta _i^m=0,\underset{t\mathrm{}}{lim}\theta _i^m=0.$$ (13) A key aspect of the system (7) or (9) is that it is fully non-linear. All functions in the coefficients on the left-hand side of the equation, like $`\stackrel{~}{ϵ}(t,\text{x})`$, and in the sources on the right-hand side are inhomogeneous and depend on $`\zeta _i^m`$ and $`\theta _i^m`$ via a basic set of three constraint equations: $`_i\mathrm{ln}a`$ $`=`$ $`_i\mathrm{ln}H={\displaystyle \frac{\stackrel{~}{ϵ}}{1\stackrel{~}{ϵ}}}e_{1A}\zeta _i^A,`$ (14) $`_i\varphi ^A`$ $`=`$ $`{\displaystyle \frac{\sqrt{2\stackrel{~}{ϵ}}}{\kappa }}\left(\delta _B^A+{\displaystyle \frac{\stackrel{~}{ϵ}}{1\stackrel{~}{ϵ}}}e_1^Ae_{1B}\right)\zeta _i^B,`$ (15) $`𝒟_i\mathrm{\Pi }^A`$ $`=`$ $`{\displaystyle \frac{\sqrt{2\stackrel{~}{ϵ}}}{\kappa }}H\left[(1\stackrel{~}{ϵ})\theta _i^A+\left((\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\delta _B^A\stackrel{~}{ϵ}e_1^Ae_{1B}+{\displaystyle \frac{\stackrel{~}{ϵ}}{1\stackrel{~}{ϵ}}}\stackrel{~}{\eta }^Ae_{1B}\right)\zeta _i^B\right].`$ (16) Using only these three constraints one can compute the spatial derivative of all relevant quantities, keeping in mind that $`\theta _i^m=e_{mA}^{}\theta _i^A(1\stackrel{~}{ϵ})^1Z_{mn}\zeta _i^n`$. Note that in our gauge $`𝒲`$ depends on $`t`$ only and does not get any non-linear contributions. ## III General analytic solution In this section we investigate how to solve the system (7) analytically and give formal expressions for the solution. In the next sections we will investigate cases where we can determine the solution more explicitly. We start by rewriting the system (7) into a single vector equation: $$\dot{v}_{ia}(t,\text{x})+A_{ab}(t,\text{x})v_{ib}(t,\text{x})=b_{ia}(t,\text{x}),\underset{t\mathrm{}}{lim}v_{ia}=0,v_i\left(\begin{array}{c}\zeta _i^1\\ \theta _i^1\\ \zeta _i^2\\ \theta _i^2\\ \mathrm{}\end{array}\right),b_i\left(\begin{array}{c}𝒮_i^1\\ 𝒥_i^1\\ 𝒮_i^2\\ 𝒥_i^2\\ \mathrm{}\end{array}\right).$$ (17) Here the indices $`a,b,\mathrm{}`$ label the components within this $`2N`$-dimensional space (with $`N`$ the number of fields). The matrix $`A`$ can be read off from (7) and has the following form: $`A_{2m1,2m}=1`$, $`A_{2m,2n}=\mathrm{\Theta }_{mn}`$ and $`A_{2m,2n1}=\mathrm{\Xi }_{mn}/(1\stackrel{~}{ϵ})^2`$, where $`\mathrm{\Theta }_{mn}`$ is the matrix between parentheses in the second equation of (7) and $`m,n=1,2,\mathrm{},N`$. All other entries of $`A`$ are zero. Equation (17) is non-linear since the matrix $`A(t,\text{x})`$ and the the source term $`b_i(t,\text{x})`$ are inhomogeneous functions in space and depend on the $`v_i`$ through (14)–(16). It can be solved perturbatively as an infinite hierarchy of linear equations with known source terms at each order (see also formalism ). We expand the relevant quantities as $$v_i=v_i^{(1)}+v_i^{(2)}+\mathrm{},b_i=b_i^{(1)}+b_i^{(2)}+\mathrm{},A(t,\text{x})=A^{(0)}(t)+A^{(1)}(t,\text{x})+A^{(2)}(t,\text{x})+\mathrm{}$$ (18) Then the equation of motion for $`v_i^{(m)}`$ is $$\dot{v}_{ia}^{(m)}(t,\text{x})+A_{ab}^{(0)}(t)v_{ib}^{(m)}(t,\text{x})=\stackrel{~}{b}_{ia}^{(m)}(t,\text{x}),\underset{t\mathrm{}}{lim}v_{ia}^{(m)}=0,\stackrel{~}{b}_{ia}^{(m)}b_{ia}^{(m)}\underset{j=1}{\overset{m1}{}}A_{ab}^{(mj)}v_{ib}^{(j)}.$$ (19) Let us recapitulate the meaning of the various indices, to avoid confusion. The index $`i=1,2,3`$ labels the components of spatial vectors. The indices $`A,B,\mathrm{}=1,\mathrm{},N`$ label components in field space. These indices will not occur in the rest of the paper, since they have been replaced by the indices $`m,n,\mathrm{}=1,\mathrm{},N`$ that label components in field space within the special basis as defined in (4). Next, the indices $`a,b,\mathrm{}=1,\mathrm{},2N`$ label components within the $`2N`$-dimensional space consisting of both $`\zeta `$ and $`\theta `$ as defined in (17). Finally there are the labels within parentheses that denote the order in the perturbative expansion defined above. Only with the $`i`$ and $`A,B,\mathrm{}`$ is there a difference between upper and lower indices. We now show that the source term $`\stackrel{~}{b}_i^{(m)}`$ is known from the solutions for $`v_i`$ up to order $`(m1)`$. The equation for $`v_i^{(1)}`$ is linear by construction: $`A^{(0)}`$ depends only on the homogeneous background quantities, and the only x dependence in $`b_i^{(1)}`$ is in the $`\mathrm{e}^{\mathrm{i}\text{k}\text{x}}`$, for the rest it depends on homogeneous background quantities. All of these are in the end functions of just $`H`$, $`\varphi ^A`$, and $`\mathrm{\Pi }^A`$ via their definitions. To go beyond linear order all these background quantities are perturbed as follows ($`C`$ stands for any of the quantities to be perturbed, for example $`\stackrel{~}{ϵ}`$, $`\stackrel{~}{\eta }^{}`$, etc.): $$C(t,\text{x})=C^{(0)}(t)+C^{(1)}(t,\text{x})+\mathrm{}=C^{(0)}+^2^i(_iC)^{(1)}+\mathrm{}=C^{(0)}+\overline{C}_a^{(0)}^2^iv_{ia}^{(1)}+\mathrm{}$$ (20) where we use (14)–(16) to compute $`_iC`$ and $`\overline{C}^{(0)}`$ is some homogeneous (space-independent) vector that is the result of that calculation. Next, to compute $`C^{(2)}`$ one simply repeats this process with the vector $`\overline{C}`$, and continues in this way order by order (of course there is also a $`\overline{C}_a^{(0)}^2^iv_{ia}^{(2)}`$ term at second order, etc.). Then it is easy to see that $`\stackrel{~}{b}_i^{(m)}`$ depends only on $`v_i^{(1)},\mathrm{},v_i^{(m1)}`$, and hence is a known quantity at each order. The solution of equation (19) for $`v_i^{(m)}`$ can be written as $$v_{ia}^{(m)}(t,\text{x})=_{\mathrm{}}^tdt^{}G_{ab}(t,t^{})\stackrel{~}{b}_{ib}^{(m)}(t^{},\text{x}),$$ (21) with the Green’s function $`G_{ab}(t,t^{})`$ satisfying<sup>2</sup><sup>2</sup>2To be precise, the Green’s function is actually defined as the solution of (22) with $`\delta (tt^{})`$ on the right-hand side instead of zero. The solution is then a step function times what we call the Green’s function. This step function has been taken into account by changing the upper limit of the integral in (21) from $`\mathrm{}`$ to $`t`$. $$\frac{\mathrm{d}}{\mathrm{d}t}G_{ab}(t,t^{})+A_{ac}^{(0)}(t)G_{cb}(t,t^{})=0,\underset{tt^{}}{lim}G_{ab}(t,t^{})=\delta _{ab}.$$ (22) It is important to note that this Green’s function is homogeneous, a solution of the background equation. It has to be computed only once, and can then be used to calculate the solution at each order using the different source terms as in (21). We write explicitly for the first two orders: $`b_{ia}^{(1)}(t,\text{x})`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\dot{𝒲}(k,t)X_{am}^{(1)}(k,t)\alpha _m(\text{k})\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}+\mathrm{c}.\mathrm{c}.,`$ $`b_{ia}^{(2)}(t,\text{x})`$ $`=`$ $`\left(^2^iv_{ic}^{(1)}(t,\text{x})\right){\displaystyle \frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\dot{𝒲}(k,t)\overline{X}_{amc}^{(1)}(k,t)\alpha _m(\text{k})\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}+\mathrm{c}.\mathrm{c}.`$ (23) Comparison with (10) shows that $`X_{am}`$ is given by the following equations: $$X_{2n1,m}=\frac{\kappa }{a\sqrt{2\stackrel{~}{ϵ}}}Q_{nm}^{\mathrm{lin}},X_{2n,m}=\frac{\kappa }{a\sqrt{2\stackrel{~}{ϵ}}}\left[\dot{Q}_{nm}^{\mathrm{lin}}Q_{np}^{\mathrm{lin}}\frac{(1+\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\delta _{pn}+Z_{pn}}{1\stackrel{~}{ϵ}}\right].$$ (24) The quantity $`\overline{X}_{amc}`$ is derived from $`X_{am}`$ as in (20), but in addition it also contains the perturbation of the basis vector $`e_m`$ inside the $`\alpha _m`$. In the same way we define $`A_{ab}^{(1)}(t,\text{x})=\overline{A}_{abc}^{(0)}(t)^2^iv_{ic}^{(1)}(t,\text{x})`$. Using the solution (21), valid at each order, and the relations (11) to compute the averages, it is now straightforward to write down the general expressions for the two-point and three-point correlators of the adiabatic ($`m=1`$) component of $`\zeta ^m^2^i\zeta _i^m`$, which is the $`a=1`$ component of $`v_{ia}`$, or rather their Fourier transforms, the power spectrum and the bispectrum. Making use of the short-hand notation $$v_{am}^{(1)}(k,t)_{\mathrm{}}^tdt^{}G_{ab}(t,t^{})\dot{𝒲}(k,t^{})X_{bm}^{(1)}(k,t^{}),$$ (25) we find for the power spectrum: $$\|\zeta ^{(1)\mathrm{\hspace{0.17em}1}}(k,t)\|^2=v_{1m}^{(1)}(k,t)v_{1m}^{(1)}(k,t)+\mathrm{c}.\mathrm{c}.$$ (26) We emphasise here that the quantity $`v_{am}^{(1)}`$ defined by (25) is simply made up of the linear $`\zeta _{mn}^{\mathrm{lin}}`$, $`\theta _{mn}^{\mathrm{lin}}`$ mode functions. One needs to evaluate the Green’s function $`G_{ab}(t,t^{})`$ and to perform the integral (25) when the linear source terms $`Q_{mn}^{\mathrm{lin}}`$, $`\dot{Q}_{mn}^{\mathrm{lin}}`$ in $`X_{bm}^{(1)}`$ are known only up to horizon crossing, $`kaH`$. However, where analytic solutions are available for $`kaH`$, or in a fully numerical approach, we can dispense with the integral (25) by using the linear solution on super-horizon scales. To find the bispectrum the calculation is slightly longer. One first has to compute $`\zeta ^{(2)\mathrm{\hspace{0.17em}1}}(t,\text{x})`$. As we noted in the single-field case in sf , $`\zeta ^{(2)\mathrm{\hspace{0.17em}1}}`$ is indeterminate. To remove this ambiguity and also require that perturbations have a zero average, we define $`\stackrel{~}{\zeta }^m\zeta ^m\zeta ^m`$. Expanding $`\stackrel{~}{\zeta }^m=\stackrel{~}{\zeta }^{(1)m}+\stackrel{~}{\zeta }^{(2)m}`$, the three-point correlator will be a combination of the different permutations of $`\stackrel{~}{\zeta }^{(2)\mathrm{\hspace{0.17em}1}}(t,\text{x}_1)\stackrel{~}{\zeta }^{(1)\mathrm{\hspace{0.17em}1}}(t,\text{x}_2)\stackrel{~}{\zeta }^{(1)\mathrm{\hspace{0.17em}1}}(t,\text{x}_3)`$, and the bispectrum is its Fourier transform. The intermediate steps are given in more detail in the explicit calculation in section V; here we go directly to the end result for the bispectrum of the adiabatic component: $$\stackrel{~}{\zeta }^1(t,\text{x}_1)\stackrel{~}{\zeta }^1(t,\text{x}_2)\stackrel{~}{\zeta }^1(t,\text{x}_3)^{(2)}(\text{k}_1,\text{k}_2,\text{k}_3)=(2\pi )^3\delta ^3(\text{k}_1+\text{k}_2+\text{k}_3)\left[f(\text{k}_1,\text{k}_2)+f(\text{k}_1,\text{k}_3)+f(\text{k}_2,\text{k}_3)\right]$$ (27) with $`f(\text{k},\text{k}^{}){\displaystyle \frac{k^2+\text{k}\text{k}^{}}{\|\text{k}+\text{k}^{}\|^2}}v_{1m}^{(1)}(k,t){\displaystyle _{\mathrm{}}^t}dt^{}`$ $`G_{1a}(t,t^{})\left[\overline{X}_{amc}^{(1)}(k,t^{})\dot{𝒲}(k,t^{})\overline{A}_{abc}^{(0)}(t^{})v_{bm}^{(1)}(k,t^{})\right]`$ $`\times [v_{1n}^{(1)}(k^{},t)v_{cn}^{(1)}(k^{},t^{})+v_{1n}^{(1)}(k^{},t)v_{cn}^{(1)}(k^{},t^{})]+\mathrm{c}.\mathrm{c}.+\text{k}\text{k}^{}.`$ (28) If $`Q_{mn}^{\mathrm{lin}}`$ is real, this simplifies to $`f(\text{k},\text{k}^{})`$ $`4{\displaystyle \frac{k^2+\text{k}\text{k}^{}}{\|\text{k}+\text{k}^{}\|^2}}v_{1m}^{(1)}(k,t)v_{1n}^{(1)}(k^{},t)`$ $`\times {\displaystyle _{\mathrm{}}^t}\mathrm{d}t^{}G_{1a}(t,t^{})[\overline{X}_{amc}^{(1)}(k,t^{})\dot{𝒲}(k,t^{})\overline{A}_{abc}^{(0)}(t^{})v_{bm}^{(1)}(k,t^{})]v_{cn}^{(1)}(k^{},t^{})+\text{k}\text{k}^{}.`$ (29) This integral expression is a key result of this paper. Using our methodology, the three-point correlator with full momentum dependence has been expressed as a single time integral over quantities determined by the background model and the linear perturbations, that is, respectively the matrix $`\overline{A}_{abc}^{(0)}`$ and the solution $`Q_{mn}^{\mathrm{lin}}`$ embedded in $`\overline{X}_{amc}^{(1)}`$ (24) and $`v_{am}^{(1)}`$ (25) (in both of which the background is also implicit). Of course, one also has to find the Green’s function $`G_{ab}`$ from (22), but, like the equation for $`Q_{\mathrm{lin}B}^A`$ in formalism , this is a linear ordinary differential equation for which there is no serious impediment to finding a numerical solution, in cases where an analytic or semi-analytic solution is unknown. The integral solution (29), then, demonstrates that the calculation of the three-point correlator is straightforward and tractable. It is, in principle, similar to calculations of the power spectrum, where accurate estimates can be found from background quantities, for example, in the slow-roll approximation. Here, we only have to supplement this with the amplitudes of the linear perturbation mode functions $`Q_{mn}^{\mathrm{lin}}(k,t)`$ and the closely related Green’s function $`G_{ab}(t,t^{})`$. In section V.4, using this methodology, we provide some quantitative semi-analytic results for the bispectrum of a two-field inflation model with a quadratic potential. Before closing this section a final comment is in order. A feature of (29) which may at first sight cast doubt on its utility for quantitative calculations is its apparent dependence on the ad hoc choice for the functional form of the window function $`𝒲(k)`$. Closer inspection reveals that the second term of (29) (the $`\overline{A}`$ term) does not depend on $`𝒲(k)`$ for scales sufficiently larger than the horizon. This is evident from the fact that (25) is simply the solution to linear theory smoothed on scales larger than the horizon. Any properly normalised window function with $`𝒲(k)1`$ for scales sufficiently larger than the horizon will produce the same final answer. The $`\overline{A}`$ term represents the non-linear evolution on superhorizon scales and, as we show below, it describes an integrated effect which can lead to large non-Gaussianity. In contrast, the $`\overline{X}`$ term arises from perturbations around horizon crossing and may depend on $`𝒲`$. We find below that this term is localized around horizon crossing and that it does not give rise to observationally interesting effects. Section V.5 provides more discussion on these points. In mf2 we explicitly show that taking a step function instead of a Gaussian as window function does not change the leading-order integrated effects. ## IV Slow-roll approximation The perturbation quantity $`Q_{mn}^{\mathrm{lin}}`$ can be computed exactly numerically, or analytically within the slow-roll approximation where all slow-roll parameters are assumed to be smaller than unity. The latter was done in vantent to next-to-leading order in slow roll:<sup>3</sup><sup>3</sup>3Compared with the solution in vantent there is an extra factor of $`1/\sqrt{2}`$. It has to be introduced to take into account the difference between the classical Gaussian random numbers $`\alpha `$, which have $`\alpha \alpha ^{}=\alpha ^{}\alpha `$, and the quantum creation/annihilation operators $`\widehat{a}^{}`$, $`\widehat{a}`$, which have $`\widehat{a}^{}\widehat{a}=0`$. In sf we introduced this factor of $`1/\sqrt{2}`$ in the analogue of equation (10), which leads of course to identical results. $$Q_{mn}^{\mathrm{lin}}=\frac{1}{2\sqrt{k}}\left[E\left(\frac{c}{kR}\right)^{1+D}\right]_{mn},$$ (30) where the matrices $`D`$ and $`E`$ are defined by $$D_{mn}\stackrel{~}{ϵ}\delta _{mn}+2\stackrel{~}{ϵ}\delta _{m1}\delta _{n1}\frac{V_{mn}}{3H^2},E_{mn}(1\stackrel{~}{ϵ})\delta _{mn}+\left(2\gamma _E\mathrm{ln}2\right)D_{mn},$$ (31) with $`\gamma _E`$ Euler’s constant. Overall unitary factors that are physically irrelevant have been omitted. Using this expression the source terms are given by $`𝒮_i^m`$ $`=`$ $`{\displaystyle \frac{\kappa }{a\sqrt{2\stackrel{~}{ϵ}}}}{\displaystyle \frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\dot{𝒲}(k)Q_{mn}^{\mathrm{lin}}(k)\alpha _n(\text{k})\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}+\mathrm{c}.\mathrm{c}.,`$ $`𝒥_i^m`$ $`=`$ $`{\displaystyle \frac{\kappa }{a\sqrt{2\stackrel{~}{ϵ}}}}\left[D_{mn}(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\delta _{mn}Z_{mn}\right]{\displaystyle \frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\dot{𝒲}(k)Q_{np}^{\mathrm{lin}}(k)\alpha _p(\text{k})\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}+\mathrm{c}.\mathrm{c}.`$ (32) Now when computing $`\overline{X}_{amc}^{(1)}`$ as defined in (23), or any higher-order terms in the perturbative expansion, we would in principle have to make the background quantities in $`Q_{mn}^{\mathrm{lin}}`$ dependent on x and perturb them according to (20). However, from (30) we see that $`Q_{mn}^{\mathrm{lin}}`$ depends on x only beyond leading order in slow roll (to leading order it is just given by $`Q_{mn}^{\mathrm{lin}}=c/(2k^{3/2}R)\delta _{mn}`$). Hence in a leading-order slow-roll approximation the only non-linear parts in the source terms are the factors in front of the integrals in (32), plus the basis vector $`e_p`$ inside the $`\alpha _p`$. Within the leading-order slow-roll approximation, we now look at the two-field case, to make things a bit more explicit. In that case the matrix $`A`$ in (17) is given by $$A=\left(\begin{array}{cccc}0& 1& 0& 0\\ 0& 3& 6\stackrel{~}{\eta }^{}& 0\\ 0& 0& 0& 1\\ 0& 0& 3\chi & 3\end{array}\right),\chi \frac{V_{22}}{3H^2}+\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{}.$$ (33) The quantity $`\chi `$ is first order in slow roll. Here we used (9) and the relation, valid to leading order in slow roll (see e.g. vantent ), $$\frac{V_{1m}}{3H^2}=\frac{V_{m1}}{3H^2}=(\stackrel{~}{ϵ}\stackrel{~}{\eta }^{})\delta _{m1}\stackrel{~}{\eta }^{}\delta _{m2}.$$ (34) Using the constraints (14)–(16) we can compute the spatial derivatives that are needed to calculate $`\overline{A}`$. Some of these were given in a general form in formalism ; to first order in slow roll for the $`\theta `$ coefficients and to second order for the $`\zeta `$ coefficients we find in the two-field case, $`(_i\mathrm{ln}a)^{(1)}`$ $`=`$ $`(_i\mathrm{ln}H)^{(1)}=\stackrel{~}{ϵ}\zeta _i^{(1)\mathrm{\hspace{0.17em}1}},`$ $`(_i\mathrm{ln}\stackrel{~}{ϵ})^{(1)}`$ $`=`$ $`2\theta _i^{(1)\mathrm{\hspace{0.17em}1}}2(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\zeta _i^{(1)\mathrm{\hspace{0.17em}1}}+2\stackrel{~}{\eta }^{}\zeta _i^{(1)\mathrm{\hspace{0.17em}2}},`$ (35) $`(_i\stackrel{~}{\eta }^{})^{(1)}`$ $`=`$ $`\stackrel{~}{\eta }^{}\theta _i^{(1)\mathrm{\hspace{0.17em}1}}+(33\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\theta _i^{(1)\mathrm{\hspace{0.17em}2}}\left((\stackrel{~}{ϵ}2\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}+\stackrel{~}{\xi }^{}\right)\zeta _i^{(1)\mathrm{\hspace{0.17em}1}}+\left(3\chi +(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}(\stackrel{~}{\eta }^{})^2\right)\zeta _i^{(1)\mathrm{\hspace{0.17em}2}},`$ $`(_i\chi )^{(1)}`$ $`=`$ $`(35\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\theta _i^{(1)\mathrm{\hspace{0.17em}1}}3\stackrel{~}{\eta }^{}\theta _i^{(1)\mathrm{\hspace{0.17em}2}}\psi _1\zeta _i^{(1)\mathrm{\hspace{0.17em}1}}\left(\psi _2+(610\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{}+3\chi )\stackrel{~}{\eta }^{}\right)\zeta _i^{(1)\mathrm{\hspace{0.17em}2}},`$ introducing the two second-order slow-roll quantities $`\psi _1`$ and $`\psi _2`$ as short-hand notation: $`\psi _1`$ $``$ $`2\stackrel{~}{ϵ}\chi +(\stackrel{~}{ϵ}\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}+3(\stackrel{~}{\eta }^{})^2+\stackrel{~}{\xi }^{}+{\displaystyle \frac{\sqrt{2\stackrel{~}{ϵ}}}{\kappa }}{\displaystyle \frac{V_{221}}{3H^2}}=\mathrm{\hspace{0.25em}2}\stackrel{~}{ϵ}(\chi +2\stackrel{~}{\eta }^{})+4(\stackrel{~}{\eta }^{})^2{\displaystyle \frac{\sqrt{2\stackrel{~}{ϵ}}}{\kappa }}{\displaystyle \frac{1}{3H^2}}\left(V_{111}V_{221}\right),`$ $`\psi _2`$ $``$ $`(11\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{}3\chi )\stackrel{~}{\eta }^{}+\stackrel{~}{\xi }^{}+{\displaystyle \frac{\sqrt{2\stackrel{~}{ϵ}}}{\kappa }}{\displaystyle \frac{V_{222}}{3H^2}},`$ (36) with $`V_{mnp}e_m^Ae_n^Be_p^C𝒟_C𝒟_B_AV`$. The reason for this specific definition of $`\psi _2`$ will become clear later on. Since we have only two fields, the notation $`\stackrel{~}{\xi }^{}`$ is unambiguous. To compute $`_i\chi `$ we used that in the two-field case, because of the orthonormality of the basis vectors, $`𝒟_ie_2^A=e_1^A(e_{2B}^{}𝒟_ie_1^B)`$. All slow-roll parameters in these expressions take their homogeneous background values. From this we find that the rank-3 matrix $`\overline{A}_{abc}`$ (defined below (24)) is $$\overline{A}=\left(\begin{array}{cccc}\mathrm{𝟎}& \mathrm{𝟎}& \mathrm{𝟎}& \mathrm{𝟎}\\ \mathrm{𝟎}& \mathrm{𝟎}& 6((\stackrel{~}{ϵ}2\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}\stackrel{~}{\xi }^{},\stackrel{~}{\eta }^{},3\chi +(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}(\stackrel{~}{\eta }^{})^2,33\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})& \mathrm{𝟎}\\ \mathrm{𝟎}& \mathrm{𝟎}& \mathrm{𝟎}& \mathrm{𝟎}\\ \mathrm{𝟎}& \mathrm{𝟎}& 3(\psi _1,35\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{},\psi _2(610\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{}+3\chi )\stackrel{~}{\eta }^{},3\stackrel{~}{\eta }^{})& \mathrm{𝟎}\end{array}\right).$$ (37) In the same way we find that the matrices $`X_{am}`$ and $`\overline{X}_{amc}`$, defined in (24), are given by $`X`$ $`=`$ $`{\displaystyle \frac{\kappa }{a\sqrt{2\stackrel{~}{ϵ}}}}\left(\begin{array}{cc}1& 0\\ 0& 2\stackrel{~}{\eta }^{}\\ 0& 1\\ 0& \chi \end{array}\right){\displaystyle \frac{\mathrm{e}^t}{2k^{3/2}}},\overline{X}={\displaystyle \frac{\kappa }{a\sqrt{2\stackrel{~}{ϵ}}}}\left(\begin{array}{cc}(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{},1,\stackrel{~}{\eta }^{},0)& (\stackrel{~}{\eta }^{},0,\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{},1)\\ 2\stackrel{~}{\eta }^{}(\stackrel{~}{\eta }^{},0,\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{},1)& \text{X}_{22}\\ (\stackrel{~}{\eta }^{},0,\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{},1)& (2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{},1,\stackrel{~}{\eta }^{},0)\\ \chi (\stackrel{~}{\eta }^{},0,\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{},1)& \text{X}_{42}\end{array}\right){\displaystyle \frac{\mathrm{e}^t}{2k^{3/2}}},`$ (38) $`\text{X}_{22}`$ $``$ $`2((\stackrel{~}{ϵ}+3\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}\stackrel{~}{\xi }^{},2\stackrel{~}{\eta }^{},3\chi +(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}2(\stackrel{~}{\eta }^{})^2,33\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{}),`$ $`\text{X}_{42}`$ $``$ $`(\psi _1(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\chi ,3+5\stackrel{~}{ϵ}\stackrel{~}{\eta }^{}\chi ,\psi _2+(610\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{}+4\chi )\stackrel{~}{\eta }^{},3\stackrel{~}{\eta }^{}),`$ where we used (30), (32), and (35). Note that $`(_i\alpha _1)/\alpha _2=(_i\alpha _2)/\alpha _1=\stackrel{~}{\eta }^{}\zeta _i^1+(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\zeta _i^2+\theta _i^2`$ to leading order in slow roll. In general the Green’s function $`G_{ab}(t,t^{})`$ cannot be expressed in closed form, since the time-dependent matrix $`A`$ does not commute at different times. It can be formally represented as $$G(t,t^{})=𝒯\mathrm{exp}\left[_t^{}^tA(s)ds\right],$$ (39) where $`𝒯`$ denotes a time-ordered exponential: $$𝒯\mathrm{exp}\left[_t^{}^tA(s)𝑑s\right]\mathrm{𝟏}_t^{}^tA(s)ds+_t^{}^tds_t^{}^sds^{}A(s)A(s^{})_t^{}^tds_t^{}^sds^{}_t^{}^s^{}ds^{\prime \prime }A(s)A(s^{})A(s^{\prime \prime })+\mathrm{}.$$ (40) This formal expression is standard in quantum mechanics and quantum field theory (see e.g. P&S ) where, viewed as a perturbative expansion, the first few terms in the series are kept when the operator $`A`$ contains a small parameter. In our case, however, not all elements of $`A`$ are first order in slow roll, so that a truncation at any finite order is a bad approximation. Moreover, even if $`A`$ were first order in slow roll, one should still be careful, because the time interval in the integrations can easily be of the order of an inverse slow-roll parameter. The time-ordered exponential can be written as an ordinary exponential plus terms which contain (nested) commutators. For example, the second and third order terms in the series (40) can be written as $$_t^{}^tds_t^{}^sds^{}A(s)A(s^{})=\frac{1}{2}\left(_t^{}^tA(s)ds\right)^2+\frac{1}{2}_t^{}^tds_t^{}^tds^{}[A(s),A(s^{})]\mathrm{\Theta }(ss^{})$$ (41) and $`{\displaystyle _t^{}^t}ds{\displaystyle _t^{}^s}ds^{}{\displaystyle _t^{}^s^{}}ds^{\prime \prime }`$ $`A(s)A(s^{})A(s^{\prime \prime })={\displaystyle \frac{1}{3!}}\left({\displaystyle _t^{}^t}A(s)𝑑s\right)^3+{\displaystyle \frac{1}{2}}{\displaystyle _t^{}^t}ds{\displaystyle _t^{}^t}ds^{}{\displaystyle _t^{}^t}ds^{\prime \prime }A(s)[A(s^{}),A(s^{\prime \prime })]\mathrm{\Theta }(s^{}s^{\prime \prime })`$ $`+{\displaystyle \frac{1}{3}}{\displaystyle _t^{}^t}ds{\displaystyle _t^{}^t}ds^{}{\displaystyle _t^{}^t}ds^{\prime \prime }\left([A(s),A(s^{})]A(s^{\prime \prime })A(s^{\prime \prime })[A(s),A(s^{})]\right)\mathrm{\Theta }(ss^{})\mathrm{\Theta }(ss^{\prime \prime })`$ (42) respectively, where $`\mathrm{\Theta }`$ is the step function, and similarly for higher orders (see lam for general expressions). There are basically three ways to proceed with this expression. In the first place we can, if we are interested only in relatively short time intervals, neglect the commutator terms in the expansion of the time-ordered exponential and write it as an ordinary exponential. The commutator terms all contain a difference of slow-roll parameters at different times, as opposed to the terms of the ordinary exponential that have just a slow-roll parameter at one time. Hence, for small time intervals, the commutator terms are a slow-roll order of magnitude smaller. Then we have an exact analytic solution in closed form for the Green’s function. Secondly, as will be the case in the explicit example in the next section, we can consider examples where $`A`$ does commute with itself at different times, in which case the time-ordered exponential simplifies to an ordinary exponential exactly. Finally, we can compute the Green’s function numerically and use it in a semi-analytic calculation (remember that the Green’s function has to be computed only once). That will be worked out in a future publication, though we give some results in section V.4. ## V Explicit solution for two-field slow-roll case In this section we provide an analytic solution for the bispectrum in two-field slow-roll inflation. We assume slow roll as in section IV but in order to obtain explicit solutions in sections V.1, V.2 and V.3 we make the further assumption that all slow-roll parameters are constant. The semi-analytic results of section V.4 are not bound by this further assumption and all slow-roll parameters are calculated numerically for a quadratic model. No slow-roll parameter takes values greater than unity. However, a semi-analytic approach is feasible for the most general non-slow-roll case and will be the study of a future publication RSvTnum . ### V.1 Power spectrum We now restrict ourselves to just two fields. Moreover, we assume the background values of $`H`$ and all slow-roll parameters, including perpendicular ones, to be completely constant in time whenever they are the leading-order (in slow roll) non-zero terms in our expressions. Then we can actually solve the system explicitly, i.e. do the time integrals. We start from the equation of motion (17) for $`v_i`$ together with the definitions (33), or rather from (19) for the $`m`$th order $`v_i^{(m)}`$ in an expansion in perturbation orders. We assume everywhere that $`\chi >0`$. At first order in the perturbations this reads as $$\dot{v}_i^{(1)}+A^{(0)}v_i^{(1)}=b_i^{(1)},\underset{t\mathrm{}}{lim}v_i^{(1)}=0.$$ (43) The matrix $`A^{(0)}`$ contains just background quantities, which by assumption are constant. Hence we circumvent the issues of non-commutativity and time-ordered exponentials, and we can write down the solution immediately as $$v_i^{(1)}(t,\text{x})=\mathrm{e}^{A^{(0)}t}_{\mathrm{}}^tdt^{}\mathrm{e}^{A^{(0)}t^{}}b_i^{(1)}(t^{},\text{x}).$$ (44) (In the terminology of section III, the Green’s function is $`G(t,t^{})=\mathrm{exp}[A^{(0)}(tt^{})]`$.) The exponential can be worked out using its eigenvalues and eigenvectors. When multiplied with its inverse (at a different time) and expanded to first order in slow roll, we obtain $`\mathrm{e}^{A^{(0)}t}\mathrm{e}^{A^{(0)}t^{}}=`$ (45) $`\left(\begin{array}{cccc}1& \frac{1}{3}\left(1(y/y^{})^3\right)& \frac{2\stackrel{~}{\eta }^{}}{\chi }\left(1(1+\frac{\chi }{3})(y/y^{})^\chi +\frac{\chi }{3}(y/y^{})^{3\chi }\right)& \frac{2\stackrel{~}{\eta }^{}}{3\chi }\left(1(y/y^{})^3(1+\frac{2\chi }{3})\left((y/y^{})^\chi (y/y^{})^{3\chi }\right)\right)\\ 0& (y/y^{})^3& 2\stackrel{~}{\eta }^{}\left((y/y^{})^\chi (y/y^{})^{3\chi }\right)& \frac{2\stackrel{~}{\eta }^{}}{\chi }\left((y/y^{})^3+\frac{\chi }{3}(y/y^{})^\chi (1+\frac{\chi }{3})(y/y^{})^{3\chi }\right)\\ 0& 0& (1+\frac{\chi }{3})(y/y^{})^\chi \frac{\chi }{3}(y/y^{})^{3\chi }& \frac{1}{3}(1+\frac{2\chi }{3})\left((y/y^{})^\chi (y/y^{})^{3\chi }\right)\\ 0& 0& \chi \left((y/y^{})^\chi (y/y^{})^{3\chi }\right)& \frac{\chi }{3}(y/y^{})^\chi +(1+\frac{\chi }{3})(y/y^{})^{3\chi }\end{array}\right)`$ Obviously, it is the identity matrix if $`y^{}=y`$. For calculational simplicity here we have defined the new time variable $`y`$, as well as the relative momentum $`p`$, $$y\frac{k_{}c}{\sqrt{2}}\mathrm{e}^t=\frac{k_{}R}{\sqrt{2}}=\frac{c}{\sqrt{2}}\mathrm{e}^{\mathrm{\Delta }t_{}},p\frac{k}{k_{}}py=\frac{kR}{\sqrt{2}}=\frac{c}{\sqrt{2}}\mathrm{e}^{\mathrm{\Delta }t_k},$$ (46) with $`\mathrm{\Delta }t_{}tt_{}`$, the time since horizon crossing of a reference mode $`k_{}`$, and we have used the fact that $`k_{}\mathrm{exp}(t_{})=1`$ by definition ($`\mathrm{\Delta }t_k`$ is defined similarly for the mode $`k`$). The fixed reference mode $`k_{}`$ is most conveniently chosen to be one of the observable modes, say the one that crossed the horizon 50 e-folds before the end of inflation. From (23) and (38) and the relation $`a=kc/(pyH\sqrt{2})`$ we find that $`b_i^{(1)}`$ can be written to leading order in slow roll as $$b_i^{(1)}=\frac{\kappa }{2\sqrt{2}}\frac{H}{\sqrt{\stackrel{~}{ϵ}}}\frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\frac{1}{k^{3/2}}\mathrm{\hspace{0.17em}2}p^2y^2\mathrm{e}^{p^2y^2}\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}\left(\begin{array}{c}\alpha _1(\text{k})\\ 2\stackrel{~}{\eta }^{}\alpha _2(\text{k})\\ \alpha _2(\text{k})\\ \chi \alpha _2(\text{k})\end{array}\right)+\mathrm{c}.\mathrm{c}.$$ (47) Changing to $`y`$ as integration variable, we can then do the integral in (44) explicitly to find the solution $`v_i^{(1)}(y,\text{x})`$ $`=`$ $`{\displaystyle \frac{\kappa }{2\sqrt{2}}}{\displaystyle \frac{H}{\sqrt{\stackrel{~}{ϵ}}}}{\displaystyle \frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\frac{1}{k^{3/2}}\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}\left(\begin{array}{cccc}\alpha _1+2\frac{\stackrel{~}{\eta }^{}}{\chi }\alpha _2& 0& 2\frac{\stackrel{~}{\eta }^{}}{\chi }\alpha _2& 0\\ 0& 0& 2\stackrel{~}{\eta }^{}\alpha _2& 0\\ 0& 0& \alpha _2& 0\\ 0& 0& \chi \alpha _2& 0\end{array}\right)\left(\begin{array}{c}\mathrm{e}^{p^2y^2}\\ p^3y^3\mathrm{\Gamma }(\frac{1}{2},p^2y^2)\\ p^\chi y^\chi \mathrm{\Gamma }(1\frac{\chi }{2},p^2y^2)\\ p^{3\chi }y^{3\chi }\mathrm{\Gamma }(\frac{1}{2}+\frac{\chi }{2},p^2y^2)\end{array}\right)}+\mathrm{c}.\mathrm{c}.`$ (48) $``$ $`{\displaystyle \frac{\kappa }{2\sqrt{2}}}{\displaystyle \frac{H}{\sqrt{\stackrel{~}{ϵ}}}}{\displaystyle \frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\frac{1}{k^{3/2}}\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}B(\text{k})u(py)}+\mathrm{c}.\mathrm{c}.`$ where we have omitted the explicit k dependence of the $`\alpha `$’s and the final expression defines the matrix $`B(\text{k})`$ and vector $`u(py)`$. It is interesting to look at the time behaviour of (48) in more detail. Note that in our leading-order slow-roll approximation, differences $`\mathrm{\Delta }t`$ in the time variable defined in (3) are equal to differences in the number of e-folds. A few e-folds<sup>4</sup><sup>4</sup>4For example, 3 e-folds is good enough if $`c=3`$ and $`\chi =0.05`$, and this result depends only weakly on the values of $`c`$ and $`\chi `$. after horizon crossing the vector $`u(py)`$ in (48) can be approximated by $`(1,0,p^\chi y^\chi \mathrm{\Gamma }(1\chi /2),0)^T`$. The third entry can be approximated even further as just $`1\chi \mathrm{\Delta }t_k`$, where $`\mathrm{\Delta }t_k`$ is the number of e-folds after horizon crossing of the mode $`k`$, and the expression is valid for $`\chi \mathrm{\Delta }t_k`$ sufficiently smaller than unity, but $`\mathrm{\Delta }t_k\text{ }\begin{array}{c}>\\ \end{array}3`$.<sup>5</sup><sup>5</sup>5See the previous footnote. A logarithmic dependence on $`c`$ has been ignored here. For $`c=3`$ this term is 4 times smaller than $`\chi \mathrm{\Delta }t_k`$ when $`\mathrm{\Delta }t_k=3`$, and becomes even less important as $`\mathrm{\Delta }t_k`$ grows. With this approximation the solution (48) can be written as $$v_i^{(1)}(t,\text{x})\frac{\kappa }{2\sqrt{2}}\frac{H}{\sqrt{\stackrel{~}{ϵ}}}\frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\frac{1}{k^{3/2}}\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}\left(\begin{array}{c}\alpha _1+2\stackrel{~}{\eta }^{}\mathrm{\Delta }t_k\alpha _2\\ 2\stackrel{~}{\eta }^{}(1\chi \mathrm{\Delta }t_k)\alpha _2\\ (1\chi \mathrm{\Delta }t_k)\alpha _2\\ \chi (1\chi \mathrm{\Delta }t_k)\alpha _2\end{array}\right)+\mathrm{c}.\mathrm{c}.$$ (49) As expected (see e.g. vantent ) we find that the effectively single-field ($`\alpha _1`$) component of $`\zeta _i^1`$ reaches its constant final value right after horizon crossing, while the influence of the perpendicular field direction ($`\alpha _2`$, ‘isocurvature mode’) on $`\zeta _i^1`$ continues to grow with time on super-horizon scales. The velocities $`\theta _i^1`$ and $`\theta _i^2`$ are both suppressed by an additional slow-roll factor compared to the $`\zeta _i`$’s. In the limit of $`\mathrm{\Delta }t_k\mathrm{}`$ (or $`py0`$), where the above approximation is no longer valid, the exact result (48) leads to the limit $$v_i^{(1)}(\text{x})\frac{\kappa }{2\sqrt{2}}\frac{H}{\sqrt{\stackrel{~}{ϵ}}}\frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\frac{1}{k^{3/2}}\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}\left(\begin{array}{c}\alpha _1+2\frac{\stackrel{~}{\eta }^{}}{\chi }\alpha _2\\ 0\\ 0\\ 0\end{array}\right)+\mathrm{c}.\mathrm{c}.$$ (50) Hence the expression does not diverge as $`t`$ grows, but reaches a well-defined value, which is independent of the smoothing parameter $`c`$. Concentrating now on the adiabatic ($`e_1`$) component of $`\zeta ^2^i\zeta _i`$ we find from (48) to leading order in slow roll: $$\zeta ^{(1)\mathrm{\hspace{0.17em}1}}(t,\text{x})=\frac{\kappa }{2\sqrt{2}}\frac{H}{\sqrt{\stackrel{~}{ϵ}}}\frac{\mathrm{d}^3\text{k}}{(2\pi )^{3/2}}\frac{1}{k^{3/2}}\mathrm{e}^{\mathrm{i}\text{k}\text{x}}\left[\mathrm{e}^{p^2y^2(\mathrm{\Delta }t_{})}\alpha _1(\text{k})+2\frac{\stackrel{~}{\eta }^{}}{\chi }g(p,\chi ,\mathrm{\Delta }t_{})\alpha _2(\text{k})\right]+\mathrm{c}.\mathrm{c}.,$$ (51) where $`y`$ as a function of $`\mathrm{\Delta }t_{}`$ is given in (46), $`p=k/k_{}`$, and we have defined $$g(p,\chi ,\mathrm{\Delta }t_{})\mathrm{e}^{p^2y^2}p^\chi y^\chi \mathrm{\Gamma }(1\frac{\chi }{2},p^2y^2)1p^\chi \mathrm{e}^{\chi \mathrm{\Delta }t_{}}=1\mathrm{e}^{\chi \mathrm{\Delta }t_k},$$ (52) where the approximation is good from a few e-folds after horizon crossing. Hence the two-point correlator is given by $$\zeta ^{(1)\mathrm{\hspace{0.17em}1}}(t,\text{x})\zeta ^{(1)\mathrm{\hspace{0.17em}1}}(t,\text{x}^{})=\frac{\kappa ^2}{8}\frac{H^2}{\stackrel{~}{ϵ}}\frac{\mathrm{d}^3\text{k}}{(2\pi )^3}\frac{1}{k^3}\left[\mathrm{e}^{2p^2y^2(\mathrm{\Delta }t_{})}+4\left(\frac{\stackrel{~}{\eta }^{}}{\chi }\right)^2g^2(p,\chi ,\mathrm{\Delta }t_{})\right]\mathrm{e}^{\mathrm{i}\text{k}(\text{x}\text{x}^{})}+\mathrm{c}.\mathrm{c}.,$$ (53) or, equivalently, for the power spectrum: $$\zeta ^{(1)\mathrm{\hspace{0.17em}1}}(k,t)\zeta ^{(1)\mathrm{\hspace{0.17em}1}}(k,t)=\frac{\kappa ^2}{4}\frac{H^2}{\stackrel{~}{ϵ}}\frac{1}{k^3}\left[\mathrm{e}^{2p^2y^2(\mathrm{\Delta }t_{})}+4\left(\frac{\stackrel{~}{\eta }^{}}{\chi }\right)^2g^2(p,\chi ,\mathrm{\Delta }t_{})\right].$$ (54) Here we used (11) to take the average. Alternatively, we could have used (26) directly. From a few e-folds after horizon crossing, $`\mathrm{exp}(2p^2y^2)1`$ and $`g(p,\chi ,\mathrm{\Delta }t_{})`$ is given by the final expression in (52), so that the power spectrum is independent of the smoothing parameter $`c`$. Finally, we can compute the adiabatic spectral index using the expressions in vantent2 , where the $`U_{Pe}`$ in that paper can be read off from (54), once the transient horizon-crossing effects have disappeared, to be $`2(\stackrel{~}{\eta }^{}/\chi )g(p,\chi ,\mathrm{\Delta }t_{})e_2`$, $$n_{\mathrm{ad}}1=4\stackrel{~}{ϵ}2\stackrel{~}{\eta }^{}8\stackrel{~}{\eta }^{}\frac{\stackrel{~}{\eta }^{}}{\chi }g(p,\chi ,\mathrm{\Delta }t_{})\frac{1g(p,\chi ,\mathrm{\Delta }t_{})}{1+4\left(\frac{\stackrel{~}{\eta }^{}}{\chi }\right)^2g^2(p,\chi ,\mathrm{\Delta }t_{})}.$$ (55) ### V.2 Second-order solution At second order in the perturbations we expand all quantities in $`A`$ and $`b_i`$ as explained in (20), using (35), resulting in the expressions in (37) and (38). Remember that superscripts within parentheses denote the order in perturbation theory, while the superscripts without parentheses indicate the component of the vector within the field basis as defined in (4). The resulting equation for $`v_i^{(2)}`$ has the same structure as (43), but with a different source term: $$\dot{v}_i^{(2)}+A^{(0)}v_i^{(2)}=b_i^{(2)}\left(\begin{array}{c}0\\ 6\stackrel{~}{\eta }^{(1)}\\ 0\\ 3\chi ^{(1)}\end{array}\right)\zeta _i^{(1)\mathrm{\hspace{0.17em}2}},\underset{t\mathrm{}}{lim}v_i^{(2)}=0,$$ (56) where $`b_i^{(2)}`$ is the vector obtained by perturbing $`H`$, $`\stackrel{~}{ϵ}`$, $`\stackrel{~}{\eta }^{}`$, $`\chi `$, and $`e_1`$ and $`e_2`$ inside $`\alpha _1`$ and $`\alpha _2`$ in $`b_i^{(1)}`$ given in (47). Explicitly, the right-hand side of equation (56) is to leading order in slow roll given by $`{\displaystyle \frac{\kappa ^2}{8}}{\displaystyle \frac{H^2}{\stackrel{~}{ϵ}}}`$ $`{\displaystyle }{\displaystyle }{\displaystyle \frac{\mathrm{d}^3\text{k}\mathrm{d}^3\text{k}^{}}{(2\pi )^3}}{\displaystyle \frac{\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}{k^{\frac{3}{2}}k_{}^{}{}_{}{}^{\frac{3}{2}}}}[2p^2y^2\mathrm{e}^{p^2y^2}\left(\begin{array}{cccc}(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\alpha _1(\text{k})\stackrel{~}{\eta }^{}\alpha _2(\text{k})& \alpha _1(\text{k})& \stackrel{~}{\eta }^{}\alpha _1(\text{k})(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\alpha _2(\text{k})& \alpha _2(\text{k})\\ & & & \\ \stackrel{~}{\eta }^{}\alpha _1(\text{k})+(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\alpha _2(\text{k})& \alpha _2(\text{k})& (\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\alpha _1(\text{k})\stackrel{~}{\eta }^{}\alpha _2(\text{k})& \alpha _1(\text{k})\\ & & & \end{array}\right)`$ $`3\left(\begin{array}{c}0\\ 0\\ 1\\ 0\end{array}\right)^TB(\text{k})u(py)\left(\begin{array}{cccc}0& 0& 0& 0\\ 2(\stackrel{~}{ϵ}2\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}+2\stackrel{~}{\xi }^{}& 2\stackrel{~}{\eta }^{}& 6\chi 2(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}+2(\stackrel{~}{\eta }^{})^2& 6+6\stackrel{~}{ϵ}2\stackrel{~}{\eta }^{}\\ 0& 0& 0& 0\\ \psi _1& 35\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{}& \psi _2(610\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{}+3\chi )\stackrel{~}{\eta }^{}& 3\stackrel{~}{\eta }^{}\end{array}\right)]`$ $`\times (B(\text{k}^{})u(qy)\mathrm{e}^{\mathrm{i}\text{k}^{}\text{x}}+\mathrm{c}.\mathrm{c}.)+\mathrm{c}.\mathrm{c}.`$ $`{\displaystyle \frac{\kappa ^2}{8}}`$ $`{\displaystyle \frac{H^2}{\stackrel{~}{ϵ}}}{\displaystyle }{\displaystyle }{\displaystyle \frac{\mathrm{d}^3\text{k}\mathrm{d}^3\text{k}^{}}{(2\pi )^3}}{\displaystyle \frac{\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}{k^{\frac{3}{2}}k_{}^{}{}_{}{}^{\frac{3}{2}}}}[2p^2y^2\mathrm{e}^{p^2y^2}\stackrel{~}{B}(\text{k})3(0,0,1,0)B(\text{k})u(py)F](B(\text{k}^{})u(qy)\mathrm{e}^{\mathrm{i}\text{k}^{}\text{x}}+\mathrm{c}.\mathrm{c}.)+\mathrm{c}.\mathrm{c}.`$ (57) where we have defined $`qk^{}/k_{}`$, as well as the matrices $`\stackrel{~}{B}(\text{k})`$ and $`F`$ in the last expression (the matrix $`B`$ and vector $`u`$ were defined in (48)). The entries indicated by an asterisk in the matrix $`\stackrel{~}{B}`$ are not given explicitly here, but can be read off from (38); they do not contribute to $`\zeta _i^{(1)\mathrm{\hspace{0.17em}1}}`$ and $`\zeta _i^{(1)\mathrm{\hspace{0.17em}2}}`$ to leading order in slow roll, because they are one order higher than the corresponding entries in the first and third row, after cancellations in the final result have been taken into account. The solution for $`v_i^{(2)}(t,\text{x})`$ is now given by the same expression (44) as $`v_i^{(1)}(t,\text{x})`$, if one replaces $`b_i^{(1)}`$ in that expression by (V.2), though actually calculating the time integral to obtain a completely explicit expression is clearly more difficult. To get all the time-dependent terms together, it is useful to change from the matrix notation used above to a component notation, as defined in (17). We define the indices $`a,b,c,d,e,f`$ running from 1 to 4 to label the components in the 4-dimensional $`\{\zeta _i^1,\theta _i^1,\zeta _i^2,\theta _i^2\}`$ space. Moreover, we rewrite the matrix in (45) as $$\mathrm{e}^{A^{(0)}t}\mathrm{e}^{A^{(0)}t^{}}=K_{abc}w_c(y,y^{}),$$ (58) $$K\left(\begin{array}{cccc}(1,0,0,0)& \frac{1}{3}(1,1,0,0)& \frac{2\stackrel{~}{\eta }^{}}{\chi }(1,0,1\frac{\chi }{3},\frac{\chi }{3})& \frac{2\stackrel{~}{\eta }^{}}{3\chi }(1,1,1\frac{2\chi }{3},1+\frac{2\chi }{3})\\ \mathrm{𝟎}& (0,1,0,0)& 2\stackrel{~}{\eta }^{}(0,0,1,1)& \frac{2\stackrel{~}{\eta }^{}}{\chi }(0,1,\frac{\chi }{3},1\frac{\chi }{3})\\ \mathrm{𝟎}& \mathrm{𝟎}& (0,0,1+\frac{\chi }{3},\frac{\chi }{3})& \frac{1}{3}(0,0,1+\frac{2\chi }{3},1\frac{2\chi }{3})\\ \mathrm{𝟎}& \mathrm{𝟎}& \chi (0,0,1,1)& (0,0,\frac{\chi }{3},1+\frac{\chi }{3})\end{array}\right),w(y,y^{})\left(\begin{array}{c}1\\ (y/y^{})^3\\ (y/y^{})^\chi \\ (y/y^{})^{3\chi }\end{array}\right),$$ which defines the rank-3 matrix $`K`$ and the vector $`w(y,y^{})`$. Then the solution for $`v_i^{(2)}`$ can be written as $`v_{ia}^{(2)}(y,\text{x})`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{8}}{\displaystyle \frac{H^2}{\stackrel{~}{ϵ}}}{\displaystyle }{\displaystyle }{\displaystyle \frac{\mathrm{d}^3\text{k}\mathrm{d}^3\text{k}^{}}{(2\pi )^3}}{\displaystyle \frac{\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}}{k^{\frac{3}{2}}k_{}^{}{}_{}{}^{\frac{3}{2}}}}K_{abc}(B_{de}(\text{k}^{})\mathrm{e}^{\mathrm{i}\text{k}^{}\text{x}}+\mathrm{c}.\mathrm{c}.)`$ $`\times [\stackrel{~}{B}_{bd}(\text{k}){\displaystyle _y^{\mathrm{}}}\mathrm{d}y^{}\mathrm{\hspace{0.17em}2}p^2y^{}\mathrm{e}^{p^2y_{}^{}{}_{}{}^{2}}w_c(y,y^{})u_e(qy^{})3B_{3f}(\text{k})F_{bd}{\displaystyle _y^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}y^{}}{y^{}}}w_c(y,y^{})u_e(qy^{})u_f(py^{})]+\mathrm{c}.\mathrm{c}.`$ (To be precise, $`c`$, $`e`$, and $`f`$ are actually indices in a completely different space than the other indices, but it is also 4-dimensional.) It is useful to consider the contributions of the two integral terms within the square brackets separately, since they have a different origin (cf. (29)). The first term is the variation of the stochastic source, represented in (29) by the $`\overline{X}`$ term, and because of the window function the integral only picks up a contribution around horizon crossing (although this contribution is time-dependent even later on, because of the dependence on $`y`$, not just $`y^{}`$, of the Green’s function). The second term is the variation of the coefficients in the equation of motion, represented in (29) by the $`\overline{A}`$ term, which is an integrated effect up to the end of inflation, and is not present in single-field inflation. The leading-order coefficients in front of the first term are first order in slow roll, while the ones in front of the second term are second order<sup>6</sup><sup>6</sup>6There are some entries in the product $`F_{bd}B_{de}`$ that are first order in slow roll, but these exactly cancel when (V.2) is worked out explicitly, so that the non-vanishing leading-order coefficients are second order in slow roll., however, this can be more than compensated by the larger integration interval. In principle there are 80 different integrals here: 16 from the first term and 64 from the second one. Some of the integrals can be done analytically, but most have to be studied numerically. However, of those 64 from the second term the only integrals that matter are those that are secular, i.e. continue to grow with time (up to a time of order $`\mathrm{\Delta }t\chi ^1`$), since these will be, roughly speaking, a slow-roll order of magnitude larger at the end of inflation than the other integrals. Although we studied all integrals more carefully, one can easily get an idea of which integrals in the second term will be secular by looking at the behaviour of the integrand for $`y^{}0`$ (i.e. $`t\mathrm{}`$): only the components with $`e`$ and $`f`$ either 1 or 3 are secular, since these are close to $`y_{}^{}{}_{}{}^{1}`$ in that limit. (Actually the components with $`e`$ and $`f`$ either 2 or 4 are zero in this slow-roll approximation, as can be seen from (48)). A slightly more careful analysis shows that, roughly speaking, the $`c=2`$ and $`c=4`$ components of those terms will be a factor $`\chi `$ smaller (one gets a $`3^1`$ instead of a $`\chi ^1`$ when integrating). Given that $`B_{31}=0`$, $`f`$ cannot be equal to 1, so in the end one expects the 4 integrals in the second term with $`c`$ and $`e`$ both either 1 or 3 and $`f=3`$ to be dominant, and that is confirmed by a careful numerical study. We denote the $`(c,e,f)=(1,1,3)`$, $`(1,3,3)`$, $`(3,3,3)`$, and $`(3,1,3)`$ integrals in the second term within the square brackets of (V.2) by $`I_1(p,q,\chi ,\mathrm{\Delta }t_{})`$, $`I_2(p,q,\chi ,\mathrm{\Delta }t_{})`$, $`I_3(p,q,\chi ,\mathrm{\Delta }t_{})`$, and $`I_4(p,q,\chi ,\mathrm{\Delta }t_{})`$, respectively. Regarding the 16 integrals of the first term the following can be said. Because of the $`y^3`$ factor in front of the $`c=2,4`$ integrals, which cannot be completely canceled by factors coming from the integral, these terms will become negligible after just a few e-folds after horizon crossing. Of the remaining integrals those with $`e=2`$ and $`e=4`$ are zero. Hence there are also only 4 distinct integrals here that have to be considered: those with $`c=1,3`$ and $`e=1,3`$ Again, these simple estimates are confirmed by careful numerical study of the integrals. We denote the $`(c,e)=(1,1)`$, $`(1,3)`$, $`(3,3)`$, and $`(3,1)`$ integrals in the first term within the square brackets of (V.2) by $`J_1(p,q,\mathrm{\Delta }t_{})`$, $`J_2(p,q,\chi ,\mathrm{\Delta }t_{})`$, $`J_3(p,q,\chi ,\mathrm{\Delta }t_{})`$, and $`J_4(p,q,\chi ,\mathrm{\Delta }t_{})`$, respectively. Let us now investigate these 8 integrals. Half of them, viz. $`I_3`$, $`I_4`$, $`J_3`$, and $`J_4`$, are zero in the limit of $`t\mathrm{}`$, but decrease slowly enough with time that they should not be neglected at the end of inflation. Three of the integrals can be done analytically: $`J_1(p,q,\mathrm{\Delta }t_{})`$ $`=`$ $`{\displaystyle _{y(\mathrm{\Delta }t_{})}^{\mathrm{}}}dy^{}\mathrm{\hspace{0.17em}2}p^2y^{}\mathrm{e}^{(p^2+q^2)y_{}^{}{}_{}{}^{2}}={\displaystyle \frac{p^2}{p^2+q^2}}\mathrm{e}^{(p^2+q^2)y^2},`$ $`J_3(p,q,\chi ,\mathrm{\Delta }t_{})`$ $`=`$ $`q^\chi y^\chi (\mathrm{\Delta }t_{}){\displaystyle _{y(\mathrm{\Delta }t_{})}^{\mathrm{}}}dy^{}\mathrm{\hspace{0.17em}2}p^2y^{}\mathrm{e}^{p^2y_{}^{}{}_{}{}^{2}}\mathrm{\Gamma }(1{\displaystyle \frac{\chi }{2}},q^2y_{}^{}{}_{}{}^{2})`$ $`=`$ $`{\displaystyle \frac{q^2}{p^2+q^2}}(p^2+q^2)^{\chi /2}y^\chi \mathrm{\Gamma }(1{\displaystyle \frac{\chi }{2}},(p^2+q^2)y^2)+q^\chi y^\chi \mathrm{e}^{p^2y^2}\mathrm{\Gamma }(1{\displaystyle \frac{\chi }{2}},q^2y^2),`$ $`J_4(p,q,\chi ,\mathrm{\Delta }t_{})`$ $`=`$ $`y^\chi (\mathrm{\Delta }t_{}){\displaystyle _{y(\mathrm{\Delta }t_{})}^{\mathrm{}}}dy^{}\mathrm{\hspace{0.17em}2}p^2y_{}^{}{}_{}{}^{1\chi }\mathrm{e}^{(p^2+q^2)y_{}^{}{}_{}{}^{2}}={\displaystyle \frac{p^2}{p^2+q^2}}(p^2+q^2)^{\chi /2}y^\chi \mathrm{\Gamma }(1{\displaystyle \frac{\chi }{2}},(p^2+q^2)y^2),`$ where $`y(\mathrm{\Delta }t_{})`$ is given in (46). It is also interesting to look at the behaviour of the integrals in the limits of $`p0`$ ($`k0`$) and $`q0`$ ($`k^{}0`$). For $`p0`$ all integrals are zero. For $`q0`$ only $`I_1`$, $`I_4`$, $`J_1`$, and $`J_4`$ are non-zero. The integrals $`J_1`$ and $`J_4`$ are given above, but in this limit also $`I_1`$ and $`I_4`$ can be computed analytically (the expression for $`I_4`$ is only valid from a few e-folds after horizon crossing, i.e. for $`py1`$): $`I_1(p,0,\chi ,\mathrm{\Delta }t_{})`$ $`=`$ $`{\displaystyle _{y(\mathrm{\Delta }t_{})}^{\mathrm{}}}dy^{}p^\chi y_{}^{}{}_{}{}^{1+\chi }\mathrm{\Gamma }(1{\displaystyle \frac{\chi }{2}},p^2y_{}^{}{}_{}{}^{2})={\displaystyle \frac{1}{\chi }}\left(\mathrm{e}^{p^2y^2}p^\chi y^\chi \mathrm{\Gamma }(1{\displaystyle \frac{\chi }{2}},p^2y^2)\right)={\displaystyle \frac{1}{\chi }}g(p,\chi ,\mathrm{\Delta }t_{}),`$ $`I_4(p,0,\chi ,\mathrm{\Delta }t_{})`$ $`=`$ $`p^\chi y^\chi (\mathrm{\Delta }t_{}){\displaystyle _{y(\mathrm{\Delta }t_{})}^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}y^{}}{y^{}}}\mathrm{\Gamma }(1{\displaystyle \frac{\chi }{2}},p^2y_{}^{}{}_{}{}^{2})p^\chi y^\chi \mathrm{ln}(py)\mathrm{\Gamma }\left(1{\displaystyle \frac{\chi }{2}}\right)(\mathrm{\Delta }t_{}\mathrm{ln}p)p^\chi \mathrm{e}^{\chi \mathrm{\Delta }t_{}},`$ (61) where $`g(p,\chi ,\mathrm{\Delta }t_{})`$ is defined in (52). Using the results discussed in the text above equation (49), one sees that from a few e-folds after horizon crossing both start growing linearly with $`\mathrm{\Delta }t_k`$ ($`=\mathrm{\Delta }t_{}\mathrm{ln}p`$), although finally the limit $`1/\chi `$ is reached for $`I_1`$, while $`I_4`$ goes to zero. For $`q>0`$ the four $`I`$-integrals have to be evaluated numerically; the resuls are plotted in figure 1 as a function of $`q/p`$ for $`\mathrm{\Delta }t_{}=50`$, for various values of the parameters $`\chi `$ and $`p`$. Although we will be using the exact numerical results for all integrals when plotting the three-point correlator, one can get an approximation by neglecting the $`\chi /2`$ inside the gamma function in $`u_3(py^{})`$ (see (V.2) and (48)). Then the $`I`$-integrals can be done analytically, with the results $`I_1(p,q,\chi ,\mathrm{\Delta }t_{})`$ $`{\displaystyle \frac{1}{2}}p^\chi (p^2+q^2)^{\chi /2}\mathrm{\Gamma }({\displaystyle \frac{\chi }{2}},(p^2+q^2)y^2),`$ $`I_2(p,q,\chi ,\mathrm{\Delta }t_{})`$ $`{\displaystyle \frac{1}{2}}p^\chi q^\chi (p^2+q^2)^\chi \mathrm{\Gamma }(\chi ,(p^2+q^2)y^2),`$ $`I_3(p,q,\chi ,\mathrm{\Delta }t_{})`$ $`{\displaystyle \frac{1}{2}}p^\chi q^\chi (p^2+q^2)^{\chi /2}y^\chi \mathrm{\Gamma }({\displaystyle \frac{\chi }{2}},(p^2+q^2)y^2),`$ $`I_4(p,q,\chi ,\mathrm{\Delta }t_{})`$ $`{\displaystyle \frac{1}{2}}p^\chi y^\chi \mathrm{\Gamma }(0,(p^2+q^2)y^2),`$ (62) so that we can make the following estimates: $`I_1`$ $`\mathrm{\Delta }t_k\left(1\frac{1}{2}\chi \mathrm{\Delta }t_k+\frac{1}{6}(\chi \mathrm{\Delta }t_k)^2\right),`$ $`I_2`$ $`\mathrm{\Delta }t_k\left(1\chi \mathrm{\Delta }t_k+\frac{2}{3}(\chi \mathrm{\Delta }t_k)^2\right),`$ $`I_3`$ $`\mathrm{\Delta }t_k\left(1\frac{3}{2}\chi \mathrm{\Delta }t_k+\frac{7}{6}(\chi \mathrm{\Delta }t_k)^2\right),`$ $`I_4`$ $`\mathrm{\Delta }t_k\left(1\chi \mathrm{\Delta }t_k+\frac{1}{2}(\chi \mathrm{\Delta }t_k)^2\right),`$ (63) for $`\chi ^1\mathrm{\Delta }t_k`$, and $`I_1`$ $`\chi ^1,`$ $`I_2`$ $`\frac{1}{2}\chi ^1,`$ $`I_3`$ $`0,`$ $`I_4`$ $`0,`$ (64) for $`\chi ^1\mathrm{\Delta }t_k`$. As a rough approximation, they can be taken independent of $`q`$ for reasonable ranges, say up to $`q/p100`$. For $`I_2`$ and $`I_3`$ this range has a lower limit as well: $`100^1\text{ }\begin{array}{c}<\\ \end{array}q/p\text{ }\begin{array}{c}<\\ \end{array}100`$; they are zero for $`q=0`$. Note that these secular $`I`$-integrals typically give a result which is of the order of an inverse slow-roll parameter. For the $`J`$-integrals the results are much smaller, since the integration interval is restricted because of the window function. The integrals $`J_1`$ and $`J_2`$ become completely independent of $`\mathrm{\Delta }t_k`$ from a few e-folds after horizon crossing of the mode $`k`$. Moreover, $`J_1`$ is independent of $`\chi `$, while $`J_2`$ has a relatively weak dependence on $`\chi `$. On the other hand, both depend strongly on $`q/p`$. They are plotted in figure 2(a). The integrals $`J_3`$ and $`J_4`$ depend strongly on both $`q/p`$ and $`\chi \mathrm{\Delta }t_k`$. In the limit $`\chi \mathrm{\Delta }t_k1`$ they become equal to $`J_2`$ and $`J_1`$, respectively, while in the opposite limit they both go to zero. They are plotted in figure 2(b). For $`q=0`$ we have an exact analytic result; for $`q=p`$ (i.e. $`k=k^{}`$) it is sometimes useful to have an analytic approximation: $`J_1`$ $`\frac{1}{2},`$ $`J_2`$ $`\frac{1}{2},`$ $`J_3`$ $`\frac{1}{2}(1\chi \mathrm{\Delta }t_k),`$ $`J_4`$ $`\frac{1}{2}(1\chi \mathrm{\Delta }t_k),`$ (65) for $`\chi ^1\mathrm{\Delta }t_k`$, and $`J_1`$ $`\frac{1}{2},`$ $`J_2`$ $`\frac{1}{2},`$ $`J_3`$ $`0,`$ $`J_4`$ $`0,`$ (66) for $`\chi ^1\mathrm{\Delta }t_k`$. Having studied all the integrals, we can now work out (V.2) explicitly. We focus on the $`a=1`$ component of $`v_{ia}`$, that is $`\zeta _i^1`$ (the adiabatic component of $`\zeta _i`$), since that is the quantity we want to compute the three-point correlator of in the end. The final result for $`\zeta _i^{(2)\mathrm{\hspace{0.17em}1}}`$ in the two-field case, in a leading-order slow-roll approximation (constant slow-roll parameters) and valid well after horizon crossing, is $`\zeta _i^{(2)\mathrm{\hspace{0.17em}1}}(t,\text{x})`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{8}}{\displaystyle \frac{H^2}{\stackrel{~}{ϵ}}}{\displaystyle \frac{\mathrm{d}^3\text{k}\mathrm{d}^3\text{k}^{}}{(2\pi )^3}\frac{1}{k^{\frac{3}{2}}k_{}^{}{}_{}{}^{\frac{3}{2}}}\mathrm{e}^{\mathrm{i}\text{k}^{}\text{x}}}`$ $`\times [(\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}\alpha _1(\text{k})+\mathrm{c}.\mathrm{c}.)\{[(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})J_1+{\displaystyle \frac{2(\stackrel{~}{\eta }^{})^2}{\chi }}(J_1J_4)]\alpha _1(\text{k}^{})`$ $`+2{\displaystyle \frac{\stackrel{~}{\eta }^{}}{\chi }}[(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})(J_1J_3)\stackrel{~}{ϵ}(J_2J_3)+{\displaystyle \frac{2(\stackrel{~}{\eta }^{})^2}{\chi }}(J_1J_2+J_3J_4){\displaystyle \frac{\chi }{2}}(J_22J_3)]\alpha _2(\text{k}^{})\}`$ $`+2{\displaystyle \frac{\stackrel{~}{\eta }^{}}{\chi }}(\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}\alpha _2(\text{k})+\mathrm{c}.\mathrm{c}.)\{[(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})(J_1J_4){\displaystyle \frac{\chi }{2}}J_1+\psi _1(I_1I_4)+\omega _1I_1]\alpha _1(\text{k}^{})`$ $`+2{\displaystyle \frac{\stackrel{~}{\eta }^{}}{\chi }}[(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})(J_1J_2+J_3J_4){\displaystyle \frac{\chi }{2}}(J_12J_2+J_3){\displaystyle \frac{\chi ^2}{4(\stackrel{~}{\eta }^{})^2}}(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{}\chi )J_2`$ $`+\psi _1(I_1I_2+I_3I_4)+\omega _1(I_1I_2)+{\displaystyle \frac{\chi }{2\stackrel{~}{\eta }^{}}}\psi _2(I_2I_3)+{\displaystyle \frac{\chi ^2}{2(\stackrel{~}{\eta }^{})^2}}\omega _2I_2]\alpha _2(\text{k}^{})\}]+\mathrm{c}.\mathrm{c}.,`$ with $`\omega _1\chi (\stackrel{~}{ϵ}+2\stackrel{~}{\eta }^{}\stackrel{~}{\xi }^{}/\stackrel{~}{\eta }^{})`$ and $`\omega _2(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})\stackrel{~}{\eta }^{}+(3\stackrel{~}{ϵ}\stackrel{~}{\eta }^{})\chi +(\stackrel{~}{\eta }^{})^2`$ defined as short-hand notation. The arguments $`(p,q,\chi ,\mathrm{\Delta }t_{})`$ of the integrals have been suppressed, but it should of course be kept in mind that that is where the time dependence resides in this expression. Note that in the single-field limit, where all terms with $`\alpha _2`$ disappear and $`\stackrel{~}{\eta }^{}=0`$, we recover the result of sf . For the three-point correlator we need to know $`\zeta ^{(2)\mathrm{\hspace{0.17em}1}}^2^i\zeta _i^{(2)\mathrm{\hspace{0.17em}1}}`$, which is given by the same expression (V.2), but with $`\mathrm{e}^{\mathrm{i}\text{k}^{}\text{x}}(\mathrm{i}k_i\mathrm{e}^{\mathrm{i}\text{k}\text{x}}\alpha _1(\text{k})+\mathrm{c}.\mathrm{c}.)`$ replaced by $$\left(\frac{k^2+\text{k}\text{k}^{}}{\|\text{k}+\text{k}^{}\|^2}\mathrm{e}^{\mathrm{i}(\text{k}+\text{k}^{})\text{x}}\alpha _1(\text{k})+\frac{k^2\text{k}\text{k}^{}}{\|\text{k}\text{k}^{}\|^2}\mathrm{e}^{\mathrm{i}(\text{k}\text{k}^{})\text{x}}\alpha _1^{}(\text{k})\right)$$ (68) and the same for $`\alpha _2(\text{k})`$. ### V.3 Bispectrum As in the single-field case, $`\zeta ^{(2)\mathrm{\hspace{0.17em}1}}`$ is indeterminate. To remove this ambiguity and also require that perturbations have a zero average, we define $`\stackrel{~}{\zeta }^m\zeta ^m\zeta ^m`$. Expanding $`\stackrel{~}{\zeta }^m=\stackrel{~}{\zeta }^{(1)m}+\stackrel{~}{\zeta }^{(2)m}`$ and switching over to Fourier space, we finally arrive at our end result for the three-point correlator (or rather the bispectrum) of the adiabatic component: $$\stackrel{~}{\zeta }^1(t,\text{x}_1)\stackrel{~}{\zeta }^1(t,\text{x}_2)\stackrel{~}{\zeta }^1(t,\text{x}_3)^{(2)}(\text{k}_1,\text{k}_2,\text{k}_3)=(2\pi )^3\delta ^3(\text{k}_1+\text{k}_2+\text{k}_3)\left[f(\text{k}_1,\text{k}_2)+f(\text{k}_1,\text{k}_3)+f(\text{k}_2,\text{k}_3)\right]$$ (69) with $`f(\text{k},\text{k}^{})`$ $``$ $`{\displaystyle \frac{\kappa ^4}{16}}{\displaystyle \frac{1}{k^3k_{}^{}{}_{}{}^{3}}}{\displaystyle \frac{H^4}{\stackrel{~}{ϵ}^2}}{\displaystyle \frac{k^2+\text{k}\text{k}^{}}{\|\text{k}+\text{k}^{}\|^2}}\{(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})J_1+{\displaystyle \frac{2(\stackrel{~}{\eta }^{})^2}{\chi }}(J_1J_4)`$ $`+4\left({\displaystyle \frac{\stackrel{~}{\eta }^{}}{\chi }}\right)^2[g(q,\chi ,\mathrm{\Delta }t_{})[(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})(J_1J_3)\stackrel{~}{ϵ}(J_2J_3)+{\displaystyle \frac{2(\stackrel{~}{\eta }^{})^2}{\chi }}(J_1J_2+J_3J_4){\displaystyle \frac{\chi }{2}}(J_22J_3)]`$ $`+g(p,\chi ,\mathrm{\Delta }t_{})[(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})(J_1J_4){\displaystyle \frac{\chi }{2}}J_1+\psi _1(I_1I_4)+\omega _1I_1]]`$ $`+16\left({\displaystyle \frac{\stackrel{~}{\eta }^{}}{\chi }}\right)^4g(p,\chi ,\mathrm{\Delta }t_{})g(q,\chi ,\mathrm{\Delta }t_{})[(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})(J_1J_2+J_3J_4){\displaystyle \frac{\chi }{2}}(J_12J_2+J_3){\displaystyle \frac{\chi ^2}{4(\stackrel{~}{\eta }^{})^2}}(\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{}\chi )J_2`$ $`+\psi _1(I_1I_2+I_3I_4)+\omega _1(I_1I_2)+{\displaystyle \frac{\chi }{2\stackrel{~}{\eta }^{}}}\psi _2(I_2I_3)+{\displaystyle \frac{\chi ^2}{2(\stackrel{~}{\eta }^{})^2}}\omega _2I_2]\}+\text{k}\text{k}^{}.`$ Again, this result is valid in the two-field case, in a leading-order slow-roll approximation (constant slow-roll parameters) and valid from a sufficient number of e-folds after horizon crossing that transient effects have disappeared. The function $`g(p,\chi ,\mathrm{\Delta }t_{})`$ is given in (52), $`\chi `$, $`\psi _1`$, and $`\psi _2`$ are defined in section IV, and $`\omega _1`$ and $`\omega _2`$ are defined below (V.2). Remember that all the integrals and the function $`g(p,\chi ,\mathrm{\Delta }t_{})`$ depend on the momenta via $`p=k/k_{}`$ and $`q=k^{}/k_{}`$ and hence are affected by the interchange of k and $`\text{k}^{}`$. In the single-field limit only the first term on the first line of (V.3) remains, which agrees exactly with sf . In the limit $`k_3k_1,k_2`$ (and hence $`\text{k}_1=\text{k}_2\text{k}`$, while we also fix $`\text{k}_{}=\text{k}`$ so that we do not need to write a subscript on $`\mathrm{\Delta }t`$), all the integrals can be performed analytically and the result is (leaving aside the overall factor of $`(2\pi )^3\delta ^3(_s\text{k}_s)`$): $`\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1^{(2)}={\displaystyle \frac{\kappa ^4}{8}}{\displaystyle \frac{1}{k^3k_3^3}}{\displaystyle \frac{H^4}{\stackrel{~}{ϵ}^2}}\left(1+4\left({\displaystyle \frac{\stackrel{~}{\eta }^{}}{\chi }}\right)^2\right)`$ $`\{(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})(1+4\left({\displaystyle \frac{\stackrel{~}{\eta }^{}}{\chi }}\right)^2(1\mathrm{e}^{\chi \mathrm{\Delta }t})^2)`$ (71) $`+4\left({\displaystyle \frac{\stackrel{~}{\eta }^{}}{\chi }}\right)^2(1\mathrm{e}^{\chi \mathrm{\Delta }t})[{\displaystyle \frac{\psi _1}{\chi }}(1(1+\chi \mathrm{\Delta }t)\mathrm{e}^{\chi \mathrm{\Delta }t})+{\displaystyle \frac{\omega _1}{\chi }}(1\mathrm{e}^{\chi \mathrm{\Delta }t})]\},`$ where the term on the first line within the curly brackets comes from the $`J`$-integrals, and the term on the second line from the $`I`$-integrals. Again, this agrees with the single-field result in the limit $`\stackrel{~}{\eta }^{}0`$. Unlike the single-field case, the multiple-field result cannot be expressed in terms of the scalar spectral index and the power spectrum only (see (54) and (55), and vantent2 for expressions for the isocurvature and mixing components). Instead of the three-point correlator itself, it is actually more useful to look at the ratio of the bispectrum to the square of the power spectrum, since that ratio is related to observables like the $`f_{\mathrm{NL}}`$ parameter (more about that later). Dividing (71) by the square of (54) (one with momentum $`k`$ and the other with $`k_3`$) and taking the limit of $`\stackrel{~}{\eta }^{}/\chi 1`$, we get for the two opposite limits of $`\chi \mathrm{\Delta }t`$ that $$\frac{\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1}{(\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1)^2}=\{\begin{array}{cc}2\left(\stackrel{~}{ϵ}+3\stackrel{~}{\eta }^{}\frac{\stackrel{~}{\xi }^{}}{\stackrel{~}{\eta }^{}}\right)+\psi _1\mathrm{\Delta }t\hfill & \text{for }\chi ^1\mathrm{\Delta }t,\hfill \\ 2\left(\stackrel{~}{ϵ}+3\stackrel{~}{\eta }^{}\frac{\stackrel{~}{\xi }^{}}{\stackrel{~}{\eta }^{}}+\frac{\psi _1}{\chi }\right)\hfill & \text{for }\chi ^1\mathrm{\Delta }t.\hfill \end{array}$$ (72) Now if we assume that $`\stackrel{~}{\eta }^{}`$ is larger than the other slow-roll parameters, the dominating term in both expressions will be the $`4(\stackrel{~}{\eta }^{})^2`$ in $`\psi _1`$ (36), so that the two expressions in (72) will go to $`4(\stackrel{~}{\eta }^{})^2\mathrm{\Delta }t`$ and $`8(\stackrel{~}{\eta }^{})^2/\chi `$, respectively. Hence, while this ratio of the bispectrum to the square of the power spectrum is first order in slow roll by naive power counting (counting $`1/\mathrm{\Delta }t`$ as a slow-roll parameter), as in the single-field case, it can be much larger for models with a relatively small $`\chi `$ and relatively large $`\stackrel{~}{\eta }^{}`$. For example, $`\stackrel{~}{\eta }^{}=0.07`$ and $`\mathrm{\Delta }t=1/\chi =50`$ would already give a ratio of more than unity, so that a value about 100 times larger than in the single-field case seems well within range for multiple-field models. This is confirmed by the full plot of (71) divided by the square of (54) as a function of $`\stackrel{~}{\eta }^{}`$ and $`\chi `$ given in figure 3. It is also interesting to see that in the cases where non-Gaussianity is large, this is caused by the $`I`$-integrals (i.e. the super-horizon integrated background effects that are absent in single-field inflation): roughly speaking it boils down to $`\stackrel{~}{ϵ}J_1`$ versus $`(\stackrel{~}{\eta }^{})^2I_1`$, which gives $`\stackrel{~}{ϵ}`$ versus the smaller of $`(\stackrel{~}{\eta }^{})^2\mathrm{\Delta }t`$ and $`(\stackrel{~}{\eta }^{})^2/\chi `$, either of which can easily be two orders of magnitude larger. In the opposite limit of $`k_1=k_2=k_3k`$ (where we again set $`k_{}=k`$ so that $`\mathrm{\Delta }t`$ is unambiguous) we do not have an exact analytic result for all integrals, but we can use the approximations (63)–(66). We find the following results for the bispectrum divided by the square of the power spectrum, in the limit of $`\stackrel{~}{\eta }^{}/\chi 1`$: $$\frac{\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1}{(\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1)^2}=\{\begin{array}{cc}\frac{3}{2}\left(\frac{1}{2\mathrm{\Delta }t}+\frac{\omega _1}{\chi }+\frac{\psi _2}{2\stackrel{~}{\eta }^{}}+\frac{1}{3}\psi _1\mathrm{\Delta }t\right)\hfill & \text{for }\chi ^1\mathrm{\Delta }t,\hfill \\ \frac{3}{2}\left(\frac{\chi }{2}+\frac{\omega _1}{\chi }+\frac{\psi _2}{2\stackrel{~}{\eta }^{}}+\frac{\psi _1}{\chi }\right)\hfill & \text{for }\chi ^1\mathrm{\Delta }t.\hfill \end{array}$$ (73) If we assume once again that $`\stackrel{~}{\eta }^{}`$ is larger than the other slow-roll parameters, the dominating term in both expressions will again be the $`4(\stackrel{~}{\eta }^{})^2`$ in $`\psi _1`$, so that the two expressions will go to $`2(\stackrel{~}{\eta }^{})^2\mathrm{\Delta }t`$ and $`6(\stackrel{~}{\eta }^{})^2/\chi `$, respectively. Finally we should check that all the limits that produce large non-Gaussianity do not produce an unacceptably large spectral index at the same time. Fortunately that is not the case: from (55) we derive, under the same limits as in (72) and (73), $$n_{\mathrm{ad}}1=\{\begin{array}{cc}4\stackrel{~}{ϵ}2\stackrel{~}{\eta }^{}\frac{8(\stackrel{~}{\eta }^{})^2\mathrm{\Delta }t}{1+4(\stackrel{~}{\eta }^{})^2(\mathrm{\Delta }t)^2}\hfill & \text{for }\chi ^1\mathrm{\Delta }t,\hfill \\ 4\stackrel{~}{ϵ}2\stackrel{~}{\eta }^{}\hfill & \text{for }\chi ^1\mathrm{\Delta }t.\hfill \end{array}$$ (74) After having discussed the various momentum limits, we finally show the full dependence on the relative magnitude of the momenta of the bispectrum divided by the square of the power spectrum in figure 4(a), where we did not use any analytic approximations for the integrals. To be precise, it is actually the bispectrum given in (69) and (V.3), without the overall $`(2\pi )^3\delta ^3(_s\text{k}_s)`$ factor (but taking into account the relation between the momenta that the $`\delta `$-function implies), divided by the sum of products of power spectra (54) with different momenta, as follows: $$\stackrel{~}{f}_{\mathrm{NL}}\frac{\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1(k_1,k_2,k_3)}{\left[\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1(k_1)\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1(k_2)+\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1(k_1)\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1(k_3)+\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1(k_2)\stackrel{~}{\zeta }^1\stackrel{~}{\zeta }^1(k_3)\right]/3}.$$ (75) This quantity can be seen as a momentum-dependent version of the $`f_{\mathrm{NL}}`$ parameter often used in the literature (see e.g. ngreview ).<sup>7</sup><sup>7</sup>7There is a difference of a factor of order unity between $`\stackrel{~}{f}_{\mathrm{NL}}`$ and $`f_{\mathrm{NL}}`$ even in the equal momentum limit, caused partly by the difference between $`\zeta `$ and the gravitational potential $`\mathrm{\Phi }`$ which was used in the original definition. We now choose $`k_{}`$ to be the mode that crossed the horizon 50 e-folds before the end of inflation (i.e. we set $`\mathrm{\Delta }t_{}=50`$). The function $`\stackrel{~}{f}_{\mathrm{NL}}`$ depends on the three scalars $`k_1,k_2,k_3`$, but we can plot it in a two-dimensional triangular domain if we fix their sum, which we do by setting $`(k_1+k_2+k_3)/k_{}=3`$. This convenient way of plotting the three-point correlator in a triangle, clearly demonstrating its symmetries, was introduced in sf , and is illustrated in figure 4(b).<sup>8</sup><sup>8</sup>8Note, however, that in sf a different normalisation factor was used. The quantities on the axes are $$\gamma 2\frac{k_2k_3}{k_1+k_2+k_3},\beta \sqrt{3}\frac{k_1k_2k_3}{k_1+k_2+k_3}.$$ (76) At the vertices of the triangle one of the three momenta is equal to zero. Lines of constant $`k_s`$ are parallel to the sides of the triangle (a different side for each $`s=1,2,3`$) and $`k_s`$ increases linearly perpendicular to these. At the side itself the corresponding momentum is equal to half the total sum, $`(k_1+k_2+k_3)/2`$. In the centre of the triangle all momenta have equal length. The plot in figure 4(a) has been made for all first-order slow-roll parameters equal to $`0.05`$, except $`\stackrel{~}{\eta }^{}=0.2`$ and $`\chi =0.01`$, and all second-order slow-roll parameters equal to $`0.003`$. We see that there is a dependence on the relative magnitude of the momenta. Though not visible in the figure, this dependence is strongest very near the vertices of the triangle, which is the limit of (71), where for this specific example the value $`9.4`$ is reached. Of course logarithmically the region near the vertices covers an infinite range of magnitudes in momentum ratios. (The fact that the result is largest in the squeezed momentum limit agrees with the findings of BCZ .) The value at the centre is $`3.7`$. Assuming that a naive extrapolation of this result at the end of inflation to the time of recombination is allowed, so that the quantity plotted is indeed comparable to the observable $`f_{\mathrm{NL}}`$, we see that this model does produce sufficient non-Gaussianity to be detectable with the Planck satellite. To compare this plot with the one for the single-field case in sf one should keep in mind that there an additional factor of $`(2\stackrel{~}{ϵ}+\stackrel{~}{\eta }^{})`$ was left out (and there are some differences in the momentum normalisation factor, but that does not change the magnitude much), so that the multiple-field result is indeed about two orders of magnitude larger. ### V.4 Comparison with semi-analytic calculation using quadratic potential Of course it may be argued that the approximate model considered here is not very realistic, with all slow-roll parameters constant with time (in particular $`\stackrel{~}{\eta }^{}`$ and $`\chi `$). We should also stress that the calculation here was made under the assumption that $`\chi >0`$, which is not true for all models. While we will study more realistic models in great detail in a future publication RSvTnum , both semi-analytically and purely numerically without any approximations, for direct comparison here we present a bispectrum calculation using the Green’s function formalism outlined in section III. We have investigated a simple two-field model with a quadratic potential $`V=\frac{1}{2}m_1^2\varphi _1^2+\frac{1}{2}m_2^2\varphi _2^2`$ with $`m_1=110^5\kappa ^1`$ (the overall mass magnitude can be freely adjusted to fix the amplitude of the power spectrum). The analytic solution (30) is used as the linear source term in (23), the super-horizon Green’s function is then calculated from (22) and (33), and the bispectrum computed from these using (29) and (37). We find that for a mass ratio $`m_2/m_1=9`$ and initial conditions $`\varphi _1=\varphi _2=13\kappa ^1`$ we get relatively large non-Gaussianity: with all momenta equal, that is, at the centre of the triangle, the ratio of the bispectrum to the square of the power spectrum is, in the slow-roll limit, $$\stackrel{~}{f}_{\mathrm{NL}}=1.8,$$ (77) where we have taken horizon crossing to be 58 e-folds before the end of inflation. The ratio of the contribution from the $`I`$-integrals to that of the $`J`$\- integrals is $`74`$. This confirms our assertion that the integrated secular terms (the $`\overline{A}`$ term in (29)) subsequent to horizon crossing dominate the contributions to the bispectrum. We note also that the spectral index in this model is $`0.93`$, which is observationally acceptable. While the investigation of the quadratic model is preliminary at this stage, it is clear that large non-Gaussianity ($`f_{\mathrm{NL}}`$ greater than unity) can be obtained in a real multiple-field inflation model. Even though the slow-roll parameters are definitely not constant in the quadratic model, we find that the expressions in the previous subsection can be used to make a useful approximation of the numerical result. In this case, for example, to estimate $`\stackrel{~}{\eta }^{}`$ we use a representative value and adjust $`\mathrm{\Delta }t`$ to reflect the region of support where its value is significant. To compare to the quadratic potential result, here we have taken $`(\stackrel{~}{\eta }^{})^2=0.73`$ (its maximum value) and $`\mathrm{\Delta }t=1.0`$ (full width at half maximum), and use the limit $`2(\stackrel{~}{\eta }^{})^2\mathrm{\Delta }t`$ given below (73). From this we estimate the value $`1.5`$, which is relatively close to the numerical value. This seems to indicate that the constant slow-roll, analytic results of the previous subsection<sup>9</sup><sup>9</sup>9Except expression (74) for the spectral index, which is a poor approximation in this case., while in principle unrealistic, can be used to get a first estimate of the amount of non-Gaussianity even in real models with varying slow-roll parameters.<sup>10</sup><sup>10</sup>10Note added: After our investigation of this example in an earlier version of this paper, the non-Gaussianity produced by a quadratic potential was also considered by the authors of lyth alabidi . They conclude that $`f_{NL}1`$ using the ‘$`\delta N`$ formalism’ and extrapolating results from a potential with (almost) equal masses. They have not computed the general case of unequal masses and assume deviations from the (almost) equal mass case to be small. In the case of almost equal masses, which is effectively single-field, we agree that non-Gaussianity is small. However, already in vantent it was quantitatively shown that even at linear order additional leading-order effects arise in the case of unequal masses from the effective coupling between the fields caused by the bending of the trajectory in field space. Moreover, in the model we consider here, the dominant non-Gaussianity is caused by a relatively large $`\stackrel{~}{\eta }^{}`$, so that a naive slow-roll order counting is not valid, a situation where the calculation of lyth alabidi as we understand it is not applicable. ### V.5 Discussion Before we conclude, a couple of points regarding the consistency of our approach need to be discussed. The first is an inherent limitation of our method in capturing all possible sources of non-Gaussianity, since, by using the linear perturbation solutions to source our non-linear equations, we are implicitly neglecting all non-linear interactions up to horizon crossing. We believe that this accounts for the small discrepancy between the momentum dependence of our single-field three-point correlator sf and that obtained from the tree-level action calculations of maldacena when $`k_1k_2k_3`$, whereas in the limit $`k_1,k_2k_3`$ the two correlators agree exactly. We surmise that the super-horizon non-linear effects described by our method and the horizon-crossing effects we are missing are of comparable magnitude for single-field inflation in the equal momentum limit. For multiple-field inflation, however, the situation is very different. We can see this by using the quantitative results above to interpret our key integral expression for the three-point correlator (29). In the case where multiple-field effects are large (as indicated by the behaviour of $`\stackrel{~}{\eta }^{}`$) it is the perturbation of the long-wavelength evolution term in the integrand of (29) (represented by $`\overline{A}_{abc}^{(0)}`$ and absent in the single-field case) that dominates over perturbations of the stochastic source term which contains the linear perturbations (represented by $`\overline{X}_{amc}^{(1)}`$; the term that would, in principle, be influenced by these horizon-crossing effects). For example, in the case considered in figure 4(a) the contribution of the $`\overline{X}_{amc}^{(1)}`$ term, which is given by the $`J`$-integrals in (V.3), is 62 times smaller than the contribution from the $`\overline{A}_{abc}^{(0)}`$ term at the vertices of the triangle, and 177 times smaller at the centre. But would the $`\overline{X}_{amc}^{(1)}`$ term be similarly enhanced when taking into account non-linear effects at horizon crossing in the multiple-field case? This seems unlikely, given that horizon crossing is only a short transition, while the large effects of the other term are caused by a buildup over a significant time interval. Moreover, this question appears to have been answered definitively in the negative by recent work seery . Generalising maldacena for multiple-field inflation, though only up to horizon crossing, it shows that these extra contributions remain of the order of small slow-roll parameters, just as in the single-field case. In that sense, the papers maldacena ; seery are important null results which clarify that our approach focusing on non-linear super-horizon effects will indeed capture the main non-Gaussian contributions from multiple-field inflation models. The second point regards the possible influence of loop corrections to the stochastic picture for generating and evolving inflationary perturbations. It is generally accepted within the cosmological community that quantum fluctuations can be considered classical for modes which have crossed the horizon, and we explicitly make such an assumption here by using classical random fields to set up initial conditions for long wavelengths via the source terms in (10). The long-wavelength evolution is then followed by using the classical equations of motion. A natural question to be asked is whether loop corrections might play a role in the super-horizon evolution. Recently, the question was addressed in weinberg for a single inflaton field plus a number of non-interacting massless scalar fields. A theorem was proved about the growth of loop effects and it was shown that for the theories mentioned loop effects were determined at horizon crossing and were subdominant. Since the conditions of the theorem imply $`\stackrel{~}{\eta }^{}=0`$, these results are not directly applicable to the kind of models considered in this paper. However, even if loop effects were to grow with time in such models, they would still need to dominate over the classical growth that these models can exhibit in order to interfere with the classical picture for the evolution of the perturbations we have developed here. Nevertheless, a definitive answer to such matters requires further investigation. ## VI Conclusions In this paper we have investigated non-Gaussianity in multiple-field inflation using the formalism of gp2 ; formalism , emphasizing analytic calculations. That formalism is based on fully non-linear equations for long wavelengths, with stochastic source terms taking into account the short-wavelength quantum fluctuations. For analytic calculations an expansion of the relevant equations in perturbation orders is necessary. However, it is much easier to derive the perturbed equations at second order directly from the non-linear equation of motion for $`\zeta _i`$ than from perturbing the original Einstein equations. Of course, in a fully numerical investigation no expansion in perturbation orders has to be made; this will be explored in future work. We derived two main results in this paper. The first is the general solution for the bispectrum, (27) with (28) or (29). Even though this is an integral expression, it will be relatively simple to evaluate in a semi-analytic calculation and it yields the full momentum dependence. To achieve this one only needs solutions for the homogeneous background quantities in the inflation model, for the linear perturbation variable $`Q^{\mathrm{lin}}`$ around horizon crossing, and for the homogeneous Green’s function, as well as expressions for the spatial derivatives of the various coefficient functions. The latter can all be computed analytically from the constraint equations (14)–(16); for the general two-field case all relevant expressions were given explicitly in formalism and this paper. Computing the bispectrum is then just a question of performing a few time integrals. An accurate semi-analytic treatment will be the subject of a forthcoming paper RSvTnum , though we do provide the results of a slow-roll calculation for a quadratic potential here. In the present paper, however, we have emphasized an example where we could proceed purely analytically. In the second part of the paper we studied two-field slow-roll inflation, with the strong leading-order approximation that all slow-roll parameters are constant. In this case we could work out the bispectrum explicitly analytically (apart from a few integrals that had to be done numerically, although we found analytic approximations in certain limits), which is the other main result of this paper, equation (V.3). We found that in this two-field case the bispectrum can easily be two orders of magnitude larger than in the single-field case, due to the continued buildup of non-Gaussianity on super-horizon scales caused by the influence of the isocurvature mode on the adiabatic perturbation. We note that even though the presence of isocurvature perturbations during inflation is crucial, it is not mandatory that they survive afterwards. In fact they feed into the adiabatic perturbation and can disappear by the end of inflation (as in the cases studied here). On the other hand, if any isocurvature modes do persist at the end of inflation, their fate will depend on the details of reheating and further evolution. The bispectrum divided by the square of the power spectrum, which can be seen as a momentum-dependent generalisation of the $`f_{\mathrm{NL}}`$ observable, can be $`f_{\mathrm{NL}}𝒪(1)`$$`𝒪(10)`$, or even larger in extreme cases. If a straightforward extrapolation of this result at the end of inflation to the time of recombination is justified, a subject which still needs to be studied in more detail, this means that the Planck satellite, and to a lesser extent even the WMAP satellite, should be able to confirm or rule out certain classes of multiple-field inflation models. Finally we want to stress that beyond estimating the amplitude of the bispectrum, we also give its explicit momentum dependence. While this dependence is rather flat for momenta of comparable size, there is a significant difference between more extreme momentum limits. We believe this paper is a significant step towards providing quantitative and testable predictions of non-Gaussianity from multiple-field inflation. Nevertheless there are still a number of issues that remain to be investigated in more detail, and on which we are working for future publications. In the first place, we will apply our general solution for the bispectrum to more realistic inflation models, particularly those strongly motivated by fundamental theory. This will require a semi-analytic treatment, but because we are dealing with integral equations, we believe that the strong slow-roll approximations presented here will actually provide reasonable analytic estimates of the exact results. As a first step we presented here the results of the semi-analytic slow-roll treatment of an explicit two-field model with a quadratic potential. This will be investigated in more detail in mf2 ; RSvTnum , but the results confirm the fact that non-Gaussianity can be large in multiple-field inflation models, and that our analytic approximations provide a good estimate. Next, it is of course important to study the further evolution of non-Gaussianities after inflation through recombination to the present day (see ngreview ; BMR for some work in this direction). In this paper we restricted ourselves to computing only the bispectrum of the adiabatic component of $`\zeta `$, even though we have the solution for all components. In future work we will investigate isocurvature and mixed bispectra as well. Finally, we will test our results with a purely numerical implementation of our formalism, which can also be applied to non-slow-roll models where our analytic approximations fail. In this case, the real-space realisations for $`\zeta `$ that result allow for other measures of non-Gaussianity to be determined, not just the three-point (or higher) correlator. After the disappointing results for single-field inflation, primordial non-Gaussianity is now back as an important quantitative tool for confirming or ruling out multiple-field inflation models, offering an exciting new window on the early universe as observations continue to improve over the next 5–10 years. ## Acknowledgements We thank the organisers of the “The Origin of the Primordial Density Perturbation” workshop in Lancaster, UK, in March 2005 for organising such a stimulating meeting, especially regarding non-Gaussianity, and where we presented the first version of the results of section V. This research is supported by PPARC grant PP/C501676/1.
warning/0506/hep-th0506183.html	ar5iv	text	# Large 𝑁 Reductions and Holography ## I Introduction The large $`N`$ limit of gauge theories leads to a drastic reduction of dynamical degrees of freedom EK . The quantities in a gauge theory in $`D`$ dimension can be calculated from a much simpler reduced model, which is obtained by dropping off the space-time dependence of the original gauge theory. The crucial condition for this large $`N`$ reduction to take place, in the case of $`SU(N)`$ gauge group, is a homogeneous distribution of the eigenvalues of gauge fields, which preserves the $`(Z_N)^D`$ symmetry. This is essentially because the homogeneous distribution generates space-time momentum from the gauge group BHN ; Prsi ; DW ; GK . On the otherhand, the celebrated Maldacena’s duality conjecture malda states that the large $`N`$ gauge theories have dual descriptions in terms of closed strings in higher dimensions, concretely realizing the large $`N`$ gauge theory–closed string duality tHooft and holography holo at the same time. It is interesting to ask how the large $`N`$ reductions are realized in the dual closed string theory via the Maldacena duality. Recently in the Gopakumar’s program towards a precise formulation of the large $`N`$ gauge theory–closed string duality Gprogram , I studied ’t Hooft-Feynman diagrams of correlation functions in gauge theories compactified on a thermal circle, to read off the corresponding dual geometries mine . It was mentioned that the technique used for calculating the thermal correlation functions was reminiscent to that appeared in the large $`N`$ reductions. However, this aspect was not investigated in depth there. In the present article, I clarify its relation to the large $`N`$ reductions, and its relevance for finding the dual holographic realization in the Maldacena duality. Some aspects of the large $`N`$ reductions will be shown to have simple explanations in the dual closed string description. One of the motivations for this study is that the reduced models are convenient for putting on computers, and therefore clarifying the holographic dual description of the large $`N`$ reductions will lead to the test of the Maldacena duality by computer simulations. Another main motivation is that this gives a concrete way to obtain a closed string theory from the matrix model of M-theory or the type IIB matrix model, via the well studied Maldacena duality. ## II Large $`N`$ reductions in Maldacena duality In this section, I first review and extend the argument of mine for how to probe the dual geometry of the $`(Z_N)^D`$ symmetric phase by the correlation functions in gauge theories compactified on a $`D`$ dimensional torus. Then, the large $`N`$ reductions of the gauge theories are obtained as a limit where the size of the torus are taken to zero. The issue of stability of the homogeneous distribution will be discussed with the comparison with the stability of the corresponding dual geometries. As an example I take $`D=4`$ case, where the boundary description is naturally identified with some $`SU(N)`$ gauge theory. Throughout this article I will work in the planar limit tHooft . I study the case where all fields are in adjoint representation of the gauge group.<sup>1</sup><sup>1</sup>1One can also introduce fields in fundamental representation and repeat the arguments similar to the one below, but baryons may be missed from such arguments based on Feynman diagrams Baryons . One may still expect from the dual holographic descriptions similar to what is discussed in this article that the large $`N`$ reductions still take place. When there are fermions, I put periodic boundary conditions on them in all the compactified directions.<sup>2</sup><sup>2</sup>2In the thermal case, fermions obey the anti-periodic boundary condition in the Euclidean time direction. As long as the phase is in the $`(Z_N)^D`$ symmetric phase the following argument apply, but whether which phase is realized depends also on these boundary conditions. This is necessary for obtaining the reduced model which reproduces the original gauge theory results. The crucial condition for the large $`N`$ reduction to take place is that the gauge field takes the configuration $`A_\mu ={\displaystyle \frac{1}{R_\mu }}\text{diag}(\theta _\mu ^1,\mathrm{},\theta _\mu ^N)`$ (1) in an appropriate gauge, where $`\theta _\mu ^a`$ distributing homogeneously between $`[\frac{1}{2},\frac{1}{2}]`$. The square expectation value of the Wilson loops winding around cycles of $`T^4`$ are order parameters of the $`(Z_N)^4`$ symmetry. They vanish in the $`(Z_N)^D`$ symmetric phase: $`\|W\|^2=0`$ center , where $`W={\displaystyle \frac{1}{N}}\text{Tr}P\mathrm{exp}i{\displaystyle _0^{2\pi R_\mu }}A_\mu 𝑑x^\mu ,\mu =0,\mathrm{},3.`$ (2) Here $`R_\mu `$ is the compactification radius in $`\mu `$-th direction and $`P`$ denotes the path ordering. Whether the above configuration is realized or not depends on the theory, here I am interested in a class of theories where this is the case. However, see the discussions on the Eguchi-Kawai reduction below. Suppose one calculates some field theory correlator $`𝒪_1(K_{\mu 1})\mathrm{}𝒪_n(K_{\mu n})`$ of gauge invariant local trace operators $`𝒪(K_{\mu j})`$, $`\text{Tr}\mathrm{\Phi }^{I_{j1}}\mathrm{}\mathrm{\Phi }^{I_{m_jj}}(K_{\mu j})`$ for example. Here, $`K_{\mu j}`$ is an external momentum of the $`j`$-th operator which takes integer values in the unit of $`\frac{1}{R_\mu }`$, and $`\mathrm{\Phi }^I`$’s are adjoint scalars. I take the background gauge $`D_\mu 𝒜^\mu =_\mu 𝒜^\mu +i[A_\mu ,𝒜^\mu ]=0`$, with $`𝒜^\mu `$ being the fluctuating quantum part of the gauge field and the background configuration $`A_\mu `$ being (1). I quantize the theory through the BRS formalism. Then, the momenta $`\frac{n_\mu }{R_\mu }`$ always appear in the combination $`\frac{1}{R_\mu }(n_\mu \delta _{ab}+\theta _\mu ^a\theta _\mu ^b)`$. Furthermore, in the planar limit one can always associate a loop momentum $`\frac{n_\mu ^i}{R_\mu }`$ ($`i=1,\mathrm{},\mathrm{}`$ labels the loop momentum) with an index loop $`a_i`$, and they appear in a specific combination $`\frac{1}{R_\mu }(n_\mu ^i+\theta _\mu ^{a_i})`$ Prsi ; DW ; GK ; mine .<sup>3</sup><sup>3</sup>3The origin of this combination is the covariant derivative for adjoint fields. For a planar diagram with all adjoint field propagators, the number of index loop is one more than the number of momentum loops, but one index sum factors out and just gives an overall factor $`N`$ mine . In the large $`N`$ limit one can replace the index sums with the integrations: $`{\displaystyle \underset{a_1\mathrm{}a_{\mathrm{}}}{}}G({\displaystyle \frac{\theta _\mu ^i}{R_\mu }})(N{\displaystyle \underset{\mu =0}{\overset{3}{}}}R_\mu {\displaystyle _{\frac{1}{2R_\mu }}^{\frac{1}{2R_\mu }}}𝑑P_{\mu i})G(P_{\mu i})`$ (3) where $`\frac{\theta _\mu ^i}{R_\mu }`$ was replaced with $`P_{\mu i}`$ in the $`N\mathrm{}`$ limit. As one sums over the gauge indices, the sums run over the homogenous distribution of the eigenvalues of the background gauge field. Thus the sum over the gauge indices can be replaced by the integration over the dual torus. This is the essential mechanism for the large $`N`$ reductions. The integrand for the correlator is a function of $`P_{\mu i}+\frac{n_{\mu i}}{R_\mu }`$. Hence the correlator has a form $`({\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_{\mu i}=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _{\frac{1}{2R_\mu }}^{\frac{1}{2R_\mu }}}𝑑P_{\mu i})G(P_{\mu i}+{\displaystyle \frac{n_{\mu i}}{R_\mu }},K_{\mu j})`$ (4) $`=`$ $`({\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑P_{\mu i})G(P_{\mu i},K_{\mu j}).`$ Thus the full internal loop momentum integrations of the uncompactified theory had been recovered. In other words, in the large $`N`$ limit the functional forms of the field theory correlators on $`T^4`$ with the background (1) coinside with those of the uncompactified theory (with a trivial gauge field configuration) to all orders in perturbation theory. However, notice that the external momenta $`K_{\mu j}`$ still take discrete values. Therefore when one performs Fourier transformation to the position space, one obtains the sum over images of the correlation functions of the uncompactified theory: $`.G(x_j^\mu )\|_{T^4}=.{\displaystyle \underset{m_j=\mathrm{}}{\overset{\mathrm{}}{}}}G(x_j^\mu +2\pi m_jR_\mu )\|_{R^4}.`$ (5) The result (5) was recently obtained in mine in the context of the Maldacena duality. It may be worth noting that the main ingredients in the derivation of (5) had appeared in the old study of the large $`N`$ reductions Prsi . The new viewpoint brought by mine was its bulk interpretation: (5) has a simple interpretation in the corresponding dual geometry. In the Maldacena duality, the geometry of the bulk can be probed by the gauge theory correlators. Then, (5) means that the dual geometry probed by the gauge theory Feynman diagrams of the compactified theory with the background (1) is the same to that of the uncompactified theory, except for the periodic identifications in $`T^4`$ directions mine . Recall that the result for correlation functions of composite operators (5) is not a trivial consequence of a simple compactification in the gauge theory side, but the configuration (1) was crucial: If one sums over images of each field’s Feynman propagator on $`R^4`$ to obtain the propagator on $`T^4`$ (say $`\mathrm{\Phi }^I(x_1)\mathrm{\Phi }^J(x_2)`$), which is appropriate for probing a geometry corresponding to $`A_\mu =0`$ background but not the homogeneous configuration (1), one does not obtain the sum over images of the correlation functions of the composite operators. Now I identify the large $`N`$ reduction with the zero-radii limit $`R_\mu 0`$, so that in the first line of (4) the momentum summation can be truncated to $`n_{\mu i}=0`$. Then, the original $`T^4`$ momentum summations drop out, but the gauge index summations reproduce the $`R^4`$ momentum integrations. This is the essence of the perturbative “derivation” of the holographic dual description of the large $`N`$ reduction.<sup>4</sup><sup>4</sup>4“Derivation” assuming that the Maldacena duality is correct for uncompactified theory. The “derivation” may be extended to the non-perturbative one by using the Schwinger-Dyson equation EK . There are two main options for taking the zero-radii limit, corresponding to two types of the reduced models. The “quenched” reduced models BHN are essentially the models where the condition (1) is put by hand. This is actually sufficient for a purpose of calculating quantities of the uncompactified original gauge theories. By construction the dual geometry in this case is the same as that of the uncompactified theory, up to the periodic identifications in $`T^4`$ directions. To calculate quantities which is translationally invariant along the $`T^4`$ directions from the closed string side, one just needs to study translationally invariant solutions of classical equation of motions. The periodic identifications in $`T^4`$ directions, in particular $`R_\mu 0`$ limit,<sup>5</sup><sup>5</sup>5The limit is, however, slightly subtle for conformal field theories where the small volume limit can be undone by conformal transformation (or isometry in dual closed string description). It will be more appropriate to keep $`R_\mu `$ finite in such cases. This is discussed in the next section. do not matter in this case. This is the closed string dual description of the quenched large $`N`$ reduction. The fact that in the classical bulk theories one can truncate the equation of motions to the holographic radial direction,<sup>6</sup><sup>6</sup>6In quantum theories, even if one is interested in translationally invariant quantities, space-time dimensionality comes in through the loop integrals. and the classical limit of the closed string theory corresponds to the planar limit, shows a beautiful correspondence between the two descriptions. The large $`N`$ gauge theories may be said to be “classical” in this sense. So far, I have been describing how the translationally invariant quantities can be obtained from the reduced model, but the reduced model can also be used to calculate the quantities which depend on space-time coordinates. This will be explained in the next subsection. On the otherhand, in the spirit of the original reduced model of Eguchi and Kawai EK , the configuration (1) is not put by hand, but it must be realized as a dominant saddle point. Thus, whether the large $`N`$ reduction takes place or not becomes a dynamical issue. This translates via the Maldacena duality into the issue of stability of the geometry dual to the uncompactified theory upon the zero radii limit of the $`T^4`$ compactification. The dynamical stability of the homogeneous distribution (1) against the small volume limit $`R_\mu 0`$ in gauge theories is a model dependent problem. Here I just make a few remarks on some aspects of it. In the supersymmetric case, the results of cmpct for $`S^1`$ compactification may seem to suggest the stability of the configuration (1). But since here all space-time directions are compactified, the quantum fluctuations can be suppressed only by the large $`N`$ effect. Therefore a separate study is actually in order. Below, I will discuss a role of fermions with the periodic boundary conditions, for the stability of the configuration (1). If the gauge theory contains a massless elementary fermionic field, the periodic boundary conditions on it may restrict the topology of the dual geometry to be $`R_0\times T^4`$ KK . This is because if some of the circle of $`T^4`$ shrinks to zero at some distance in the holographic radial direction in the bulk, the bulk fermion which couples to the gauge theory operator containing the massless fermion cannot have the periodic boundary conditions.<sup>7</sup><sup>7</sup>7However, this restriction may not be so strong if one takes into account other space-time directions in the dual theory. See Ross for a recent interesting example where the circle in the asymptotic boundary is mixed with another circle corresponding to an internal symmetry in field theory side. As argued above, the stability of the $`R_0\times T^4`$ topology in the bulk is necessary for the stability of the configuration (1) in the limit $`R_\mu 0`$ mine .<sup>8</sup><sup>8</sup>8The bulk topology may also be probed by using the classical closed string worldsheet as a dual description for the Wilson loop expectation values Wilson . Precisely speaking, what is calculated in Wilson is a generalization of the Wilson loop including adjoint scalars. This expectation from the closed string side may heuristically be explained in the reduced models if one recalls the procedure taken here for taking the large $`N`$ limit. To see this, I first analyze a reduced model with $`SU(2)`$ gauge group and with one massless adjoint fermion, to estimate effective potential between two eigenvalues of the gauge field. In this case it is possible to integrate out the fermion red , and it is easy to see that the presence of the massless adjoint fermion introduces a repulsive potential $`\mathrm{log}L`$ for $`L0`$, where $`L`$ is a difference between the two eigenvalues. One may expect that there is a similar repulsive force between eigenvalues also in the $`SU(N)`$ reduced model. Then, recall that to obtain the reduced models from the gauge theories, I took $`N\mathrm{}`$ before taking the $`R_\mu 0`$ limit. To implement this condition starting from the reduced models, one should restrict eigenvalues of the reduced gauge fields between $`\frac{1}{2R_\mu }`$ and $`\frac{1}{2R_\mu }`$.<sup>9</sup><sup>9</sup>9This is a gauge invariant condition for the reduced models. The mutually commuting configuration (1) should emerge dynamically. It is like putting particles with short distant repulsive forces dense enough in a finite volume, so that the resulting distribution becomes uniform, i.e. the configuration (1) is realized. On the otherhand, adding a mass term to the fermion weaken the repulsive force and it disappears if the mass is sufficiently large. It is natural because if the mass is taken large enough, the fermion will eventually decouple from the system. It suggests that in the corresponding classical solution of the dual closed string descriptions, fermionic fields which couple to a gauge invariant fermionic operator with that massive fermion are excluded from a region in the space-time corresponding to the scale lower than the fermion mass scale, and they do not restrict the topology there. The complete exclusion of fermionic field from some region may require a singular geometry in the supergravity description. The arguments given here are heuristic and deserves further study. If one introduces a bosonic adjoint field $`\mathrm{\Phi }`$ to a gauge theory instead of the fermion, it means introducing another space dimension in the dual closed string side. Here I study the simple situation where $`\mathrm{\Phi }=0`$ vacuum is realized in the gauge theory. To construct a corresponding reduced model, one should take the diagonal components of $`\mathrm{\Phi }`$ to be zero by hand, much in the same spirit as in the quenched reduced models, but for the opposite type of configuration.<sup>10</sup><sup>10</sup>10One may also try to show that this configuration is dynamically preferred in the reduced model red . Then, a calculation similar to the above shows, in $`SU(2)`$ case, that the bosonic adjoint field does not change the leading repulsive potential from the fermionic field. For purely bosonic theories, in the closed string side the $`AdS`$ soliton adssoliton which is a possible vacuum state at finite $`R_\mu `$ already partially breaks (1), and one must also take into account the possibilities of various phase transitions, like to black holes, black strings GL ; T and so on, which can trigger instability of the configuration (1) upon taking the $`R_\mu 0`$ limit. As I mentioned earlier, the configuration (1) is crucial for the large $`N`$ reduction. The instability of the geometries dual to the uncompactified theories upon compactification means that the Eguchi-Kawai reduction does not take place in those cases. ### II.1 Correlators of local gauge invariant single trace operators from reduced models As has been described in the previous subsections, the large $`N`$ reduction is not merely a simple dimensional reduction in the gauge theory side, but the configuration (1) was crucial. The loop momentum integrations in the original gauge theory are recovered from the gauge index sums in the reduced model. One can also calculate correlation functions in gauge theories which depend on external momenta from reduced models, provided that the configuration (1) is realized. This is essentially because the gauge indices in the reduced models play the role of space-time momenta in the original gauge theory. In the Maldacena duality, the correlation function of single trace operators are important because they correspond to closed string amplitudes in the dual theory. However, since trace operators do not have un-contracted gauge indices, additional techniques are required to calculate correlation functions which depend on external momenta.<sup>11</sup><sup>11</sup>11If there are un-contracted gauge indices, these can be straightforwardly regarded as the external momenta in the large $`N`$ reductions Prsi ; DW ; GK . In this subsection, I explain how to calculate such correlation functions from the reduced models. I take an operator made of adjoint scalars $`\mathrm{\Phi }^I`$ ($`I`$ labels the species of the scalars) as an example, generalizations to include fermions or dynamical gauge fields are straightforward. For an operator made of $`q`$ scalars $`\text{Tr}\mathrm{\Phi }^{I_1}\mathrm{}\mathrm{\Phi }^{I_q}(k)`$ in the gauge theory, the corresponding object in the reduced model is given by $`\text{Tr}^{\mathrm{\Delta }k}\mathrm{\Phi }^{I_1}\mathrm{}\mathrm{\Phi }^{I_q}`$ (6) where I have defined the “shifted trace” $`\text{Tr}^{\mathrm{\Delta }k}`$ as $`\text{Tr}^{\mathrm{\Delta }k}AB={\displaystyle \underset{P_1,P_2}{}}A_{P_1,P_2}B_{P_2,P_1+k}.`$ (7) I have reparametrized the matrix indices in terms of $`P_\mu `$, already taking into account the $`N\mathrm{}`$ limit. See (3). $`k_\mu `$ is regarded as a shift in the matrix indices. The point is that the shift inserts the external momentum $`k_\mu `$ to the index line connected to the latter index of $`\mathrm{\Phi }^{I_q}`$. The shifted trace does not satisfy the cyclic property $`\text{Tr}AB=\text{Tr}BA`$, so the above mapping from the gauge theory to the reduced model is not unique. However this is not a problem, since for a planar diagram these cyclically shifted operators all give the same results. This is because the index loops and the incoming momenta always appear in the combination $`P_{\mu i_1}P_{\mu i_2}+k_\mu `$, where $`i_1,i_2`$ are loop indices. $`k_\mu `$ can be assigned to either $`i_1`$ or $`i_2`$, and the difference can be absorbed by a shift of the loop momentum, which is an integration variable. See FIG. 1-3. ## III Finite radii and T-duality ### III.1 T-duality in gauge theories In the previous section I identified the large $`N`$ reductions with the dimensional reductions with the non-trivial gauge configuration (1). However, since the Maldacena duality is supposed to hold for any radii, it is natural to generalize the notion of the large $`N`$ reductions to that case. In this subsection, I will explain that the equivalence between the original gauge theory and the reduced model still holds for finite radii, by appropriately generalizing the notion of the reduced model as a matrix model on a compact space, along the line of DT . As found in DT , this naturally leads to the notion of T-dual equivalence in the matrix model. The corresponding dual closed string description of this T-dual equivalence Kikkawa will be presented in the subsequent subsection. For finite radii $`R_\mu `$, the gauge theory calculation has equivalent T-dual descriptions in terms of the matrix model DT . Each eigenvalues are interpreted as positions of D-instantons (in the string units) in the T-dual language. The radii of the dual torus $`\stackrel{~}{T}^4`$ are $`\frac{\mathrm{}_s^2}{R_\mu }`$, where $`\mathrm{}_s`$ is the string length. The summation over $`n_\mu `$ in (4) corresponds to the summation over images of D-instantons on the dual torus $`\stackrel{~}{T}^4`$. To incorporate the images in the reduced models, one embeds the $`M^4`$ $`SU(N)`$ gauge groups into the diagonal blocks in $`SU(N\times M^4)`$ gauge group, where $`M`$ is a positive integer which will be taken to infinity. The matrix components of the reduced fields are subject to an identification corresponding to the $`\stackrel{~}{T}^4`$ compactification. The background gauge field configuration (1) is generalized to $`A_{\mu \stackrel{}{m}\stackrel{}{m}}={\displaystyle \frac{1}{R_\mu }}\text{diag}(\theta _\mu ^1,\mathrm{},\theta _\mu ^N)+{\displaystyle \frac{m_\mu }{R_\mu }}`$ (8) where $`m_\mu `$ is a component of a four-vector $`\stackrel{}{m}`$ which is an index for $`SU(M^4)`$. The off-diagonal components (in terms of $`SU(M^4)`$) are zero. In the matrix model on the torus, the fields in adjoint representation satisfy DT $`\mathrm{\Phi }_{\stackrel{}{\theta }_1(\stackrel{}{m}+\stackrel{}{v}),\stackrel{}{\theta }_2(\stackrel{}{n}+\stackrel{}{v})}=\mathrm{\Phi }_{\stackrel{}{\theta }_1\stackrel{}{m},\stackrel{}{\theta }_2\stackrel{}{n}}`$ (9) where I have labeled the $`SU(N)`$ gauge group indices in terms of $`\stackrel{}{\theta }`$, and $`\stackrel{}{m},\stackrel{}{n}`$ are $`SU(M^4)`$ indices. $`\stackrel{}{v}`$ is an arbitrary four-vector with integer entries, which expresses a parallel shift to an image. For simplicity, I study massless scalar fields $`\mathrm{\Phi }^I`$. Generalization to other fields is straightforward. The quadratic term of the reduced model is given by $`{\displaystyle \frac{1}{M^4}}\text{Tr}_{SU(N\times M^4)}[A_\mu ,\mathrm{\Phi }^I][A_\mu ,\mathrm{\Phi }^I]`$ $``$ $`{\displaystyle \underset{\stackrel{}{m},\stackrel{}{\theta }_1,\stackrel{}{\theta }_2}{}}\mathrm{\Phi }_{\stackrel{}{\theta }_1\stackrel{}{m},\stackrel{}{\theta }_2\stackrel{}{0}}^I\left({\displaystyle \frac{1}{R_\mu }}(\theta _{\mu 1}\theta _{\mu 2}+m_\mu )\right)^2\mathrm{\Phi }_{\stackrel{}{\theta }_2\stackrel{}{0},\stackrel{}{\theta }_1\stackrel{}{m}}^I`$ where use has been made for (9). In the Maldacena duality, one studies the coupling of gauge invariant operators to their sources. For example, in the gauge theory the trace of $`q`$ scalar fields have the coupling of the form $`{\displaystyle d^4K𝒥_{I_1\mathrm{}I_q}(K)\text{Tr}_{SU(N)}\mathrm{\Phi }^{I_1}\mathrm{}\mathrm{\Phi }^{I_q}(K)}.`$ (11) The source $`𝒥_{I_1\mathrm{}I_q}`$ is identified with the boundary value of the corresponding field in closed string side. In the reduced models, the corresponding coupling is given by $`({\displaystyle \underset{\mu }{}}{\displaystyle _{\frac{1}{2R_\mu }}^{\frac{1}{2R_\mu }}}𝑑k_\mu {\displaystyle \underset{m_\mu =\mathrm{}}{\overset{\mathrm{}}{}}})`$ $`𝒥_{I_1\mathrm{}I_q}(K)\text{Tr}_{SU(N\times M^4)}^{\mathrm{\Delta }k,\mathrm{\Delta }\stackrel{}{m}}\mathrm{\Phi }^{I_1}\mathrm{}\mathrm{\Phi }^{I_q}`$ (12) where $`K_\mu =k_\mu +\frac{m_\mu }{R_\mu }`$, and I have introduced the “$`\stackrel{}{m}`$-shifted trace” $`\text{Tr}_{SU(M^4)}^{\mathrm{\Delta }\stackrel{}{m}}`$ for $`SU(M^4)`$ indices defined by $`\text{Tr}_{SU(M^4)}^{\mathrm{\Delta }\stackrel{}{m}}AB`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{m}_1,\stackrel{}{m}_2}{}}A_{\stackrel{}{m}_1,\stackrel{}{m}_2}B_{\stackrel{}{m}_2,(\stackrel{}{m}_1+\stackrel{}{m})}`$ (13) $`=`$ $`{\displaystyle \underset{\stackrel{}{m}_1,\stackrel{}{m}_2}{}}B_{\stackrel{}{m}_2,\stackrel{}{m}_1}A_{\stackrel{}{m}_1,(\stackrel{}{m}_2+\stackrel{}{m})}.`$ The $`SU(M^4)`$ matrices $`A`$ and $`B`$ satisfy the same condition as in (9), and the the “cyclic” property, i.e. the second line of (13) follows from that. The shifted trace for the $`SU(N)`$ indices $`P`$ is defined as in (6). One can check that the reduced model with the above source term gives the same result of the original gauge theory in the diagrammatic perturbative calculations. See the example FIG. 4,5. In the dual D-instanton descriptions, the T-dual of momenta are winding modes of closed strings, and the shift $`k`$ and $`\stackrel{}{m}`$ in the trace $`\text{Tr}_{SU(N\times M^4)}^{\mathrm{\Delta }k,\mathrm{\Delta }\stackrel{}{m}}`$ corresponds to a string stretched between a D-instanton and another D-instanton shifted by $`\mathrm{}_s^2(k_\mu +\frac{m_\mu }{R_\mu })`$. Note that usually the sum over images is not taken in the reduced models. In this sense the this is a generalized of the large $`N`$ reduction. ### III.2 The dual geometry Finally, I explain in this subsection that the generalized large $`N`$ reduction has a simple description in the dual geometry.<sup>12</sup><sup>12</sup>12I thank R. Gopakumar and K. P. Yogendran for stimulating my thought on T-dual geometries at the early stage collaboration in mine . As a concrete example, I take $`𝒩=4`$ super Yang-Mills theory on $`T^4`$, which is identified as a worldvolume theory of D3-branes, at strong coupling. At strong coupling, supergravity approximation is valid and the dual geometry is $`AdS_5\times S^5`$ with the periodic identifications in the $`T^4`$ directions, and the dilaton is constant. As will be shown below, this geometry can be obtained from a multi D-instanton solution in type IIB supergravity via T-duality, where D-instantons are densely and homogeneously distributing on the dual $`\stackrel{~}{T}^4`$.<sup>13</sup><sup>13</sup>13The T-dual relation of these geometries has appeared in Kallosh . The point of this subsection is to exhibit the parallel between the dual descriptions. The dense homogeneous distribution of the D-instantons is identified as a holographic dual of the dense and homogeneous distribution of the eigenvalues (1). Thus this is a holographic description of the equivalence between the gauge theory and the generalized reduced model. The (Euclideanized) metric for the D-instantons in Einstein frame is flat: $`\stackrel{~}{g}_{\mu \nu E}=\delta _{\mu \nu }`$, $`\mu ,\nu =0,\mathrm{},9`$ DIS . The solution can be obtained by solving the following equation for dilaton $`\stackrel{~}{\varphi }`$: $`_\mu ^\mu e^{\stackrel{~}{\varphi }}=0.`$ (14) When D-instantons are densely and homogeneously distributing in the $`\stackrel{~}{T}^4`$ directions, and overlapping on a point in the transverse six dimensions, the solution is given by $`e^{\stackrel{~}{\varphi }_{\mathrm{}}+\stackrel{~}{\varphi }}=g_s\left(1+{\displaystyle \frac{c_0g_sN\mathrm{}_s^4}{r^4}}\right)`$ (15) where $`r`$ is the radial coordinate transverse to $`\stackrel{~}{T}^4`$, $`N`$ is a number of D-instantons on $`\stackrel{~}{T}^4`$ and $`g_s=e^{\stackrel{~}{\varphi }_{\mathrm{}}}`$ is the string coupling constant. $`c_0`$ is a numerical constant related to the volume of the unit five-sphere, I suppress such numerical factors hereafter. In the near horizon limit $`r0`$, the dilaton configuration becomes $`e^{\stackrel{~}{\varphi }}={\displaystyle \frac{g_sN\mathrm{}_s^4}{r^4}}`$ (16) and I obtain the $`AdS_5\times S^5`$ metric in the string frame $`\stackrel{~}{ds}_{st}^2=e^{\stackrel{~}{\varphi }/2}\stackrel{~}{ds}_E^2`$: $`\stackrel{~}{ds}_{st}^2={\displaystyle \frac{\sqrt{g_sN\mathrm{}_s^4}}{r^2}}\left[dr^2+r^2d\mathrm{\Omega }_5+d\stackrel{~}{x}_{//}^2\right]`$ (17) where $`\stackrel{~}{x}_{//}^\mu `$ is a coordinate on $`\stackrel{~}{T}^4`$ with period $`2\pi \stackrel{~}{R}_\mu `$ <sup>14</sup><sup>14</sup>14Before taking the near horizon limit $`_{\mu =0}^3\stackrel{~}{R}_\mu =\mathrm{}_s^4`$ should hold in this solution so that $`N`$ coinsides with the number of D-instantons. After the near horizon limit this restriction can be removed by the isometry of $`AdS_5`$. and $`d\mathrm{\Omega }_5`$ is the volume form of the unit five-sphere. Now I perform T-dual transformation on $`\stackrel{~}{T}^4`$. The T-dual metric is again $`AdS_5\times S^5`$: $`ds_{st}^2={\displaystyle \frac{\sqrt{g_sN\mathrm{}_s^4}}{r^2}}\left[dr^2+r^2d\mathrm{\Omega }_5\right]+{\displaystyle \frac{r^2}{\sqrt{g_sN\mathrm{}_s^4}}}dx_{//}^2`$ (18) where $`x_{//}^\mu `$ is a coordinate on $`T^4`$ with period $`2\pi R_\mu =2\pi \frac{\mathrm{}_s^2}{\stackrel{~}{R}_\mu }`$. Under the T-duality the dilaton transforms as Buscher $`\varphi =\stackrel{~}{\varphi }{\displaystyle \frac{1}{2}}\mathrm{log}det{}_{\stackrel{~}{T}^4}{}^{}\stackrel{~}{g}_{\mu \nu st}^{}=0.`$ (19) Thus one arrives at the $`AdS_5\times S^5`$ geometry with the constant dilaton ($`e^\varphi _{\mathrm{}}=g_s`$), as I have claimed. This is the holographic description of the generalized large $`N`$ reduction in the previous subsection. Notice the key role of the dense and homogeneous distribution of the D-instantons, which is dual to the dense and homogeneous of the eigenvalues of the gauge field: It gives the geometry which is T-dual to the geometry just obtained by a simple $`T^4`$ identification of the uncompactified D3-brane near horizon geometry. Recall that the T-dual relation in the supergravity classical solutions can be derived from closed string worldsheet sigma model Kikkawa ; Buscher , whereas the matrix model T-duality was motivated and “explained” by the open stirng sigma model but was shown purely within the gauge theoretical language in DT . The validity of these two descriptions may have an overlap in the Maldacena’s large $`N`$ and the near horizon limit, as long as the conjecture is correct.<sup>15</sup><sup>15</sup>15Practically, one needs to be able to handle either the stringy corrections or the strongly coupled gauge theory. Note that although one obtains a smeared solution from D3-brane solution even for finite $`N`$, when the number of the D3-brane is small the gauge theory description is not rigorously related to this geometry. The ’t Hooft-Maldacena limit provides the correspondence between the gauge theory and the closed string theory, and the large $`N`$ limit of the homogeneous distribution of the eigenvalues of the gauge field, which is dual to the dense and homogeneous distribution of the D-instantons, provides the effective smearing of the multi-D-instanton solution. Then, the T-dual equivalence of two geometries can be interpreted as a holographic dual description of the matrix model T-dual equivalence between the gauge theory and the generalized reduced model studied in the previous subsection. ## IV Summary and Discussions In this article, I have presented the holographic dual descriptions of the large $`N`$ reductions in the Maldacena duality. This will be useful for deepening the understanding of both sides. The equivalence between the reduced model and the original gauge theory can be interpreted as a limit of the compactification with the homogeneous distribution of the eigenvalues of the gauge field. It was shown how this equivalence is reflected in the dual bulk geometry through the correlation functions of the local gauge invariant single trace operators. Since the Maldacena duality holds even for finite radii, it is natural to generalize the equivalence relation to that case. This was achieved by using the description of the matrix model on a compact space introduced in DT . This description naturally contains the notion of T-duality. I pointed out that for finite radii the T-dual equivalence of two supergravity solutions are the holographic dual description of the T-dual equivalence between the gauge theory and the generalized reduced model. The crucial condition for the large $`N`$ reduction is the homogeneous distribution of the eigenvalues of the gauge field (1). In the quenched reduced models this condition is forced by hand, whereas in the Eguchi-Kawai reduction the stability is a dynamical issue. The stability of the homogeneous distribution should reflect the stability of the supergravity solution dual to the uncompactified gauge theory upon compactification on $`T^4`$. I pointed out an interesting possible role of fermions obeying the periodic boundary conditions in the $`T^4`$ directions. I also presented a new technique for calculating position dependent correlation functions of gauge invariant single trace operators in gauge theories from the reduced models. Despite the evidences from the past studies, the Maldacena duality still remains as a conjecture. The holographic dual of the large $`N`$ reductions established in this article will be useful for the quantitative tests of the Maldacena duality. Reduced models are suitable for studying the non-perturbative effects. In the Maldacena duality, it is also expected that the classical closed string descriptions capture the non-perturbative effects of the dual gauge theories. It will be interesting to study further how non-perturbative effects in the reduced models reflect themselves in the dual closed string descriptions. The large $`N`$ reductions also provide an advantage for computer simulations NJ . As shown in this article, the large $`N`$ reductions have more direct correspondence with the Maldacena duality compared with the lattice gauge theory, at least at present. If a computer simulation of a reduced model supports the dynamical stability of the configuration (1), it suggests that there is a dual closed string solution which is stable against the limit $`R_\mu 0`$. Then one can further calculate quantities in closed string theory using the reduced model. I hope that the holographic dual descriptions of the large $`N`$ reductions described in this article will lead to the investigation of the Maldacena duality by computer simulations of reduced models. I think relating the Matrix model of M-theory BFSS or the IIB matrix model IIBM to the Maldacena duality<sup>16</sup><sup>16</sup>16I am aware that many researchers have resorted the idea of relating the matrix models with Maldacena duality from different directions. via the large $`N`$ reductions discussed in this article is the most concrete way to study how closed strings emerge from these models, especially taking into account the recent developments in the understanding of the Maldacena duality Gprogram ; mine ; spin . * * * I am grateful to my colleagues in the present and past institutions from whom I have learned a lot about the ingredients appeared in this article. I am also grateful to my colleagues in HRI for their continuous warm support. I sincerely appreciate generous supports for our research from the people in India.
warning/0506/quant-ph0506053.html	ar5iv	text	# Superballistic Diffusion of Entanglement in Disordered Spin Chains ## .1 Introduction Entanglement can be viewed as a physical resource that has no analog in classical information theory. As such, entanglement plays an important role in many quantum information tasks, such as quantum cryptography, teleportation and quantum algorithms. The development of protocols for the distribution of entanglement is an important problem in quantum information processing. Recently several methods have been proposed for accomplishing the related problem of quantum state transfer using spin chains Bose (2003); Subrahmanyam (2004); Christandl et al. (2004); Albanese et al. (2004); Christandl et al. (2005); Osborne and Linden (2004); Haselgrove (2004); Verstraete et al. (2004a, b); Plenio and Semiao (2005); Burgarth and Bose (2004, 2005); Bose et al. (2004); Kay and Ericsson (2005). An important characteristic in all of these models is that the transfer occurs with minimal control or intervention in the dynamics of the spin chain. Besides achieving perfect quantum state transfer, another important task that can be examined is that of distributing entanglement between a central node and many distributed nodes within a processor. Such a capability may be useful for the generation of multi-qubit entangled states (with the appropriate post-distillation). In any case, gaining a better understanding of how one can use spin chain dynamics to spread out (rather than transmit) entanglement is of fundamental interest, although this has not seemed to have been addressed much in the literature to date. In this letter we consider the second order spatial moment of entanglement in a spin chain, with a single excitation, and show that a finite disordered region can lead to superballistic growth in the the second order spatial moment. The entanglement properties of various spin chains have been examined in a number of recent works Latorre et al. (2004); Jin and Korepin (2004a, b); Korepin (2004); Its et al. (2005); Fan et al. (2004). It has been shown that, for a particle in a lattice with a finite disordered region, a finite period of superballistic growth can occur in the spatial variance of the wave function Hufnagel et al. (2001). This can be applied to XXZ spinchains, which have nearest neighbor interactions similar to the tight binding approximation used for a 1-dimensional lattice. Unlike the spatial variance, however, the growth in the second order spatial moment of concurrence in a chain with a finite central disordered region is not bounded from above by the similar growth experienced in an ordered chain. The disordered region leads to sustained superballistic spatial growth of the concurrence and provides an efficient mechanism for distributing entanglement along spinchains. ## .2 Concurrence One measure of bipartite entanglement is concurrence Hill and Wootters (1997), which is a monotonic function of entanglement of formation. For a mixed state, the concurrence between two qubits, $`i`$ and $`j`$, is defined to be $`C_{ij}=\mathrm{max}\{\lambda _1\lambda _2\lambda _3\lambda _4,0\}`$, where $`\lambda _n`$ is the square root of the $`n^{th}`$ eigenvalue of $`\rho \stackrel{~}{\rho }`$ in descending order. Here $`\stackrel{~}{\rho }=(\sigma _y\sigma _y)\rho ^{}(\sigma _y\sigma _y)`$, where $`\rho ^{}`$ is the complex conjugate of $`\rho `$. For a single pure excitation of a $`N`$spin 1/2 chain, the wave function is $$\|\psi =\alpha _1\|1+\mathrm{}+\alpha _N\|N,$$ (1) where $`\|k`$ is one of $`N`$ basis states of the $`N`$spin chain where the $`k^{th}`$ site is in the excited state and all other sites are in the ground state. Taking any two spins within this chain $`i`$, $`j`$, and tracing over the rest yields the following density matrix: $$\rho _{ij}=\left(\begin{array}{cccc}\mu & 0& 0& 0\\ 0& \|\alpha _i\|^2& \alpha _i\alpha _j^{}& 0\\ 0& \alpha _j\alpha _i^{}& \|\alpha _j\|^2& 0\\ 0& 0& 0& 0\end{array}\right),$$ where $`\mu =1\|\alpha _i\|^2\|\alpha _j\|^2`$. The eigenvalues of $`\rho _{ij}\stackrel{~}{\rho }_{ij}`$ are $`\{4\|\alpha _i\|^2\|\alpha _j\|^2,0,0,0\}`$, so $`\lambda _n=\{2\|\alpha _i\|\|\alpha _j\|,0,0,0\}`$. Inserting these expressions for $`\lambda _n`$ into the previous equation for concurrence gives $`C_{ij}=2\|\alpha _i\|\|\alpha _j\|`$. ## .3 Distribution of Entanglement One way to study the time evolution of the spatial distribution of entanglement between the origin site and other sites in the spin chain is to look at the second order spatial moment of concurrence about the point of origin, $$M(t)=\underset{x0}{}\|x^2C_{0x}(t)\|,$$ where $`C_{0x}(t)`$ is the concurrence at time $`t`$, between the original site of the excitation, and the site at position $`x`$. $`C_{0x}(t)`$ can be replaced with the earlier expression for the concurrence when only one excitation is acted upon by a spin preserving Hamiltonian. $`M(t)`$ $`=`$ $`{\displaystyle \underset{x0}{}}\|x^2(2\|\alpha _0(t)\|\|\alpha _x(t)\|)\|,`$ (2) $`=`$ $`2\|\alpha _0(t)\|W(t),`$ (3) where $$W(t)\underset{x0}{}x^2\|\alpha _x(t)\|.$$ (4) ## .4 Spin Dynamics: Ordered Chain A general Hamiltonian for a spin chain with nearest neighbour interactions and uniform couplings can be written as: $`H={\displaystyle \underset{i}{}}(J_x^{i,i+1}\sigma _x^i\sigma _x^{i+1}+J_y^{i,i+1}\sigma _y^i\sigma _y^{i+1}+J_z^{i,i+1}\sigma _z^i\sigma _z^{i+1})`$ In such a spin chain, with the additional constraint that couplings obey $`J_x^{i,i+1}=J_y^{i,i+1}=\mathrm{\Gamma }`$, for all sites $`i`$ in the chain, then the overall $`z`$component of spin is conserved. Here we will take $`\mathrm{\Gamma }=1`$, which will give general results up to a rescalling of $`t`$. For a spin chain of $`N`$ sites, a basis can be formed from $`N`$ vectors, $`\|1\mathrm{}\|N`$, with $`\|k`$ corresponding to a chain with a single excitation at site $`k`$. Any Hamiltonian, if of the above form in this basis, corresponds to a symmetric tridiagonal matrix. Further, if the coupling parameters $`J_z^i=0`$, then the main diagonal will also vanish, giving $`H_{ij}`$ $`=`$ $`\delta (i,j1)+\delta (i,j+1),`$ $`H`$ $`=`$ $`L+R,`$ where $`L_{ij}=\delta (i,j1)`$, $`R_{ij}=\delta (i,j+1)`$. Using this expression for $`H`$, the evolution operator, $`U(t)`$ can be found. $`U(t)`$ $`=`$ $`e^{iHt}=e^{i(L+R)t},`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}(it)^n(L+R)^n,`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}(it)^n{\displaystyle \underset{m=0}{\overset{n}{}}}{\displaystyle \frac{n!}{(nm)!m!}}L^mR^{nm},`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(it)^n{\displaystyle \underset{m=0}{\overset{n}{}}}{\displaystyle \frac{1}{(nm)!m!}}D_{n2m},`$ where $`D_k=R^k`$, and $`R^1=L`$. The only non-zero entries in $`D_k`$ are along the $`k^{th}`$ sub-diagonal, which are all 1. The family of matrices, $`\{D_k\}_{\mathrm{}}^+\mathrm{}`$, form a basis for U(t) and one can set, $`U(t)`$ $`=`$ $`{\displaystyle \underset{x=\mathrm{}}{\overset{\mathrm{}}{}}}c_x(t)D_x,`$ $`c_x(t)`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(it)^{2l+x}{\displaystyle \frac{1}{(l+x)!l!}},\text{ where }2l=nx\text{,}`$ $`=`$ $`(i)^xt^x{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(1)^l(t)^{2l}{\displaystyle \frac{1}{(l+x)!l!}},`$ $`=`$ $`(i)^xJ_x(2t),`$ where $`J_k`$ is the $`k^{th}`$ order Bessel function of the first kind. ### .4.1 Upper bound on W(t) In the case of the ordered chain we have not been able to obtain exact analytic expressions for (4), but instead we can find bounds for $`W(t)`$. Using (4), and the approximation, $`J_x(z)\{\begin{array}{c}0\text{ for }z>\|x\|,\\ \sqrt{\frac{2}{\pi z}}\mathrm{cos}(z\frac{x\pi }{2}\frac{\pi }{4}),\text{ for }z\|x\|\end{array}`$ (7) an upper bound for $`W(t)`$ can be found, $`W(t)`$ $`=`$ $`{\displaystyle \underset{x=\mathrm{}}{\overset{+\mathrm{}}{}}}x^2\|\alpha _x(t)\|={\displaystyle \underset{x=\mathrm{}}{\overset{+\mathrm{}}{}}}x^2\|J_x(2t)\|,`$ $`=`$ $`{\displaystyle \underset{x=\mathrm{}}{\overset{+\mathrm{}}{}}}\|x^2J_x(2t)\|.`$ This can be simplified using the identity $`xJ_x(a)=\frac{a}{2}(J_{x1}(a)+J_{x+1}(a))`$ Gradshteyn and Ryzhik (1980), to obtain $`W(t)`$ $`=`$ $`{\displaystyle \underset{x=\mathrm{}}{\overset{+\mathrm{}}{}}}\|t^2(J_{x2}(a)+2J_x(a)+J_{x+2}(a))`$ $`+t(J_{x1}(a)J_{x+1}(a))\|,`$ $`=`$ $`(4t^2+2t){\displaystyle \underset{x=\mathrm{}}{\overset{+\mathrm{}}{}}}\|J_x(2t)\|,`$ $``$ $`(4t^2+2t)\sqrt{{\displaystyle \frac{1}{\pi t}}}(4t),`$ $``$ $`{\displaystyle \frac{16}{\sqrt{\pi }}}t^{\frac{5}{2}}\text{ for }t1.`$ ### .4.2 Lower bound on W(t) A lower bound for $`W(t)`$ can also be found by observing $`\|J_{2k}(2t)\|`$ $``$ $`J_{2k}(2t),`$ $`{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}(2k)^2\|J_{2k}(2t)\|`$ $``$ $`{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}(2k)^2J_{2k}(2t),`$ $`{\displaystyle \underset{x=\mathrm{}}{\overset{\mathrm{}}{}}}x^2\|J_x(2t)\|`$ $``$ $`{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}(2k)^2\|J_{2k}(2t)\|,`$ $``$ $`{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}(2k)^2J_{2k}(2t).`$ Using the identities $`_{k=\mathrm{}}^{\mathrm{}}(2k)^2J_{2k}(a)=\frac{a^2}{2}`$ Gradshteyn and Ryzhik (1980), and (4), this can be rewritten, $`W(t)2t^2.`$ Since $`2t^2W(t)\frac{16}{\sqrt{\pi }}t^{\frac{5}{2}}`$ for $`t1`$, then $`W(t)`$ grows as $`t^\lambda `$, with $`2\lambda 2.5`$. ### .4.3 Approximation for M(t) Using the previous approximation (7), which can be rewritten as, $`J_x(z)\sqrt{\frac{2}{\pi z}}(\mathrm{cos}(z\frac{\pi }{4})\mathrm{cos}(\frac{x\pi }{2})+\mathrm{sin}(z\frac{\pi }{4})\mathrm{sin}(\frac{x\pi }{2}))`$, an expression for $`W(t)`$ can be found. $`W(t)`$ $`=`$ $`{\displaystyle \underset{x=\mathrm{}}{\overset{\mathrm{}}{}}}x^2\|J_x(2t)\|,`$ $``$ $`{\displaystyle \underset{x=2t}{\overset{2t}{}}}x^2\sqrt{{\displaystyle \frac{1}{\pi t}}}\|\mathrm{cos}(2t{\displaystyle \frac{x\pi }{2}}{\displaystyle \frac{\pi }{4}})\|,`$ By taking the average of $`\|\mathrm{cos}(2t\frac{x\pi }{2}\frac{\pi }{4})\|`$ over a time period from $`t=t^{}`$ to $`t=t^{}+\pi `$, the time averaged value, $`\frac{2}{\pi }`$, is obtained. Using this as an approximation for the cosine term, the expression for $`W(t)`$ reduces to: $`W(t)`$ $``$ $`{\displaystyle \frac{4}{\pi }}\sqrt{{\displaystyle \frac{1}{\pi t}}}{\displaystyle \underset{x=0}{\overset{2t}{}}}x^2,`$ $`=`$ $`{\displaystyle \frac{4}{\pi }}\sqrt{{\displaystyle \frac{1}{\pi t}}}{\displaystyle \frac{2t(2t+1)(4t+1)}{6}},`$ $``$ $`{\displaystyle \frac{32}{3\pi ^{\frac{3}{2}}}}t^{\frac{5}{2}},\text{ as }t\mathrm{}`$ This can be used in the expression for $`M(t)`$, $`M(t)`$ $`=`$ $`2\|\alpha _0(t)\|W(t),`$ $``$ $`2\sqrt{{\displaystyle \frac{1}{\pi t}}}{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{32t^{\frac{5}{2}}+24t^{\frac{3}{2}}+4t^{\frac{1}{2}}}{3\pi ^{\frac{3}{2}}}},`$ $``$ $`{\displaystyle \frac{128}{3\pi ^3}}t^2,\text{ as }t\mathrm{}.`$ ## .5 Spin Dynamics: Disordered Chain It is well known that disorder, which in this case we take to be random $`J_z`$ coupling strengths, or random magnetic fields in the $`z`$direction at each site, leads to spatial localisation of the wave function Anderson (1958). In a spin chain with a disordered region centred on the initial site of excitation, the wave function will be partially localized, and so the amplitude at the initial site will fall off at a lower rate than in the case of an ordered chain. We now model the emission of an excitation from the centre site out through a disordered region and into (on either side), a semi-infinite ordered spin chain. To consider this we first note that the dynamics of an ordered semi-infinite spin chain, with an initial single excitation at the first site, can be solved in a similar way to the infinite spin chain. For an ordered semi-infinite spin chain we have $$U(t)=e^{i(L+R)t}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}(it)^n(L+R)^n.$$ (8) For a semi-infinite chain, however, $`[L,R]0`$. $`(LR)_{i,j}=\delta _{i,j}`$, but $`(RL)_{i,j}=\delta _{i,j}(1\delta _{i,1})`$. If the initial excitation is at the start of the chain, then the resulting wave, after time $`t`$, will be the first column of $`U(t)`$. Taking the first column matrix elements of $`U(t)`$ to be the wave launched from site $`x=0`$, as $`U(t)_{x,0}`$, and we can obtain, $`U(t)_{x,0}`$ $`=`$ $`{\displaystyle \underset{n=x}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(it)^n}{n!}}\left(\left(\begin{array}{c}n\\ \frac{nx}{2}\end{array}\right)\left(\begin{array}{c}n\\ \frac{nx2}{2}\end{array}\right)\right),`$ $`=`$ $`(i)^x(J_x(2t)+J_{x+2}(2t)),`$ resulting in a wave function with elements $`\alpha _x(t)`$ (1), $`\alpha _x(t)`$ $`=`$ $`(i)^x(J_x(2t)+J_{x+2}(2t)),`$ $`=`$ $`(i)^x{\displaystyle \frac{(x+1)}{t}}J_{x+1}(2t).`$ We now suppose that the random region, extending from $`x=L`$ to $`x=+L`$, emits an excitation at $`+L`$, into the ordered semi-infinite region with amplitude $`f(t)`$. Shifting coordinates to this interface by setting $`x^{}=xL`$, and then (for convenience) dropping the prime, we can examine the wave propagation into the ordered semi-infinite chain (now $`x1`$) to be: $`\alpha _{x+L}(t)`$ $`=`$ $`{\displaystyle _0^t}f(t^{})(i)^{x1}({\displaystyle \frac{x}{tt^{}}})J_x(2(tt^{}))𝑑t^{}.`$ We now break the emission processes into two time bins. For $`tt_1`$: $`\alpha _{x+L}(t)`$ $`=`$ $`{\displaystyle _0^{t_1}}f(t^{})(i)^{x1}({\displaystyle \frac{x}{tt^{}}})J_x(2(tt^{}))𝑑t^{}`$ $`+`$ $`{\displaystyle _{t_1}^t}f(t^{})(i)^{x1}({\displaystyle \frac{x}{tt^{}}})J_x(2(tt^{}))𝑑t^{},`$ $`=`$ $`(i)^{x1}({\displaystyle \frac{x}{t\tau }})J_x(2(t\tau )){\displaystyle _0^{t_1}}f(t^{})𝑑t^{}`$ $`+`$ $`{\displaystyle _{t_1}^t}f(t^{})(i)^{x1}({\displaystyle \frac{x}{tt^{}}})J_x(2(tt^{}))𝑑t^{},`$ $`=`$ $`(i)^{x1}({\displaystyle \frac{x}{t\tau }})J_x(2(t\tau ))\beta `$ $`+`$ $`{\displaystyle _{t_1}^t}f(t^{})(i)^{x1}({\displaystyle \frac{x}{tt^{}}})J_x(2(tt^{}))𝑑t^{}.`$ For large $`t_1`$, we will neglect the second term in (LABEL:long), which is equivalent to letting $`f(t)=\beta \delta (t\tau )`$. This corresponds to a simple model for the emission in which an excitation is emited, with amplitude $`\beta `$, at time $`t=\tau t_1`$. For $`t\tau `$, there is no emission and $`\alpha _{x+L}(t)=0`$, while for $`t>\tau `$: $`\alpha _{x+L}(t)`$ $`=`$ $`\beta (i)^{x1}({\displaystyle \frac{x}{t\tau }})J_x(2(t\tau )).`$ (11) The second order moment of concurrence can be split into two parts, $`M_o(t)`$ and $`M_d(t)`$. $`M_o(t)`$ is the contribution to $`M(t)`$ from the ordered region, and $`M_d(t)`$ is the contribution from the disordered region. For $`t>\tau `$: $`M_o(t)`$ $`=`$ $`4\|\alpha _0(t)\|{\displaystyle \underset{x=1}{\overset{\mathrm{}}{}}}(x+L)^2{\displaystyle \frac{x}{t\tau }}\|\beta J_x(2(t\tau ))\|,`$ $`=`$ $`2{\displaystyle \frac{\|\beta \|}{T^{\frac{3}{2}}\sqrt{\pi }}}(\|\mathrm{cos}(2T{\displaystyle \frac{\pi }{4}})\|{\displaystyle \underset{x=1}{\overset{T}{}}}(2x+L)^22x`$ $`+\|\mathrm{sin}(2T{\displaystyle \frac{\pi }{4}})\|{\displaystyle \underset{x=1}{\overset{T}{}}}(2x+L1)^2(2x1)),`$ where $`T=t\tau `$. Since the average value of $`\|\mathrm{sin}(z)\|=\|\mathrm{cos}(z)\|=\frac{2}{\pi }`$, the above can be approximated by $`W_o(t)`$ $`=`$ $`{\displaystyle \frac{4\|\beta \|}{T^{\frac{3}{2}}\pi ^{\frac{3}{2}}}}{\displaystyle \underset{x=1}{\overset{2T}{}}}(x+L)^2x,`$ $`=`$ $`{\displaystyle \frac{2\|\beta \|}{T^{\frac{3}{2}}\pi ^{\frac{3}{2}}}}(8T^4+({\displaystyle \frac{32L}{3}}8)T^3+(4L^28L+3)T^2`$ $`+(4L22L^2)T(18+{\displaystyle \frac{16L}{3}}+2L^2)),`$ $``$ $`{\displaystyle \frac{16\|\beta \|}{\pi ^{\frac{3}{2}}}}(t\tau )^{\frac{5}{2}},\text{ for }t0\text{.}`$ From (3), to obtain $`M_o(t)`$, we require the amplitude of the wave at the origin or initial site. Since the probability of the excitation entering the semi-infinite chain to the right and left is $`\|\beta \|^2`$, the probability of the excitation remaining within the disordered region is $`12\|\beta \|^2`$. If the probability at the initial site is proportional to this, then the amplitude of the wavefunction at the initial origin will be $`\gamma \sqrt{12\|\beta \|^2}`$, with $`\|\gamma \|1`$. $`M_o(t)`$ $`=`$ $`{\displaystyle \frac{4\|\beta \gamma \sqrt{12\|\beta \|^2}\|}{T^{\frac{3}{2}}\pi ^{\frac{3}{2}}}}(8T^4+({\displaystyle \frac{32L}{3}}8)T^3`$ $`+(4L^28L+3)T^2+(4L22L^2)T`$ $`(18+{\displaystyle \frac{16L}{3}}+2L^2)),`$ $``$ $`{\displaystyle \frac{32\|\beta \gamma \sqrt{12\|\beta \|^2}\|}{\pi ^{\frac{3}{2}}}}(t\tau )^{\frac{5}{2}},\text{ for }t0\text{.}`$ There is only a finite disordered range, so $`W_d(t)`$ is bounded from above by $`2L^2`$. This in turn leads to an upper limit on $`M_d(t)`$ of $`4\|\gamma \sqrt{12\|\beta \|^2}\|L^2`$. Since $`M_ot^{\frac{5}{2}}`$ and $`M_d2\beta L^2`$, $$M(t)M_o(t)\frac{32\|\beta \gamma \sqrt{12\|\beta \|^2}\|}{\pi ^{\frac{3}{2}}}(t\tau )^{\frac{5}{2}},\text{ for }t0.$$ (15) In summary, we have examined the emission of a single excitation from a central site on a spin chain in the two cases where the chain is ordered and where the origin site is surrounded by a finite region of disorder. We have found that the spatial expansion of the two-site entanglement, between the origin site and other sites on the chain, is much more rapid in the latter case where disorder is present, than in the ordered case. Indeed the second order moment of the spatial extent of the concurrence expands super-ballistically. This effect could be used to rapidly distribute entanglement throughout regions of a spin chain. If, rather than distribute entanglement throughout a chain, one wished to spatially separate an entangled state one might instead appeal to a recent scheme outlined in Kay and Ericsson (2005). J. Fitzsimons acknowledges support from the Embark Initiative while J. Twamley acknowledges support from Science Foundation Ireland.
warning/0506/math0506588.html	ar5iv	text	# Homogenization of a diffusion process in a rarefied binary structure ## 1 Introduction Diffusion occurs naturally and is important in many industrial and geophysical problems, particularly in oil recovery, earth pollution, phase transition, chemical and nuclear processes. When one comes to a rational study of binary structures, a crucial point lies in the interaction between the microscopic and macroscopic levels and particularly the way the former influences the latter. Once the distribution is assumed to be $`\epsilon `$-periodic, this kind of study can be accomplished by the homogenization theory. The present study reveals the basic mechanism which governs diffusion in both phases of such a binary structure, formed by an ambiental connected phase surrounding a periodical suspension of small particles. For simplicity, the particles are considered here to be spheres of radius $`r_\epsilon <<\epsilon `$, that is $`\underset{\epsilon 0}{lim}{\displaystyle \frac{r_\epsilon }{\epsilon }}=0`$. We balance this assumption, which obviously means that the suspension has vanishing volume, by imposing the total mass of the suspension to be always of unity order. This simplified structure permits the accurate establishment of the macroscopic equations by means of a multiple scale method of the homogenization theory adapted for fine-scale substructures. It allows to have a general view on the specific macroscopic effects which arise in every possible case. As we use the non-dimensional framework, the discussion is made in fact with respect to only two parameters: $`r_\epsilon `$ and $`b_\epsilon `$, the latter standing for the ratio of suspension/ambiental phase diffusivities. As the diffusivities of the two components can differ by orders of magnitude, the interfacial conditions play an important role. It happens that the following cases have different treatments: $`r_\epsilon <<\epsilon ^3`$, $`\epsilon ^3<<r_\epsilon <<\epsilon `$ and $`r_\epsilon =𝒪(\epsilon ^3)`$. To give a flavor of what may be considered as an appropriate choice of the relative scales, we refer to the pioneering work where the appearance of an extra term in the limit procedure is responsible for a change in the nature of the mathematical problem and is linked to a critical size of the inclusions. Later showed how this could be generalized to the $`N`$-dimensional case for non linear operators satisfying classical properties of polynomial growth and coercivity. Since then, the notion of non local effects has been developed in a way that is closer to the present point of view in , and . In dealing with our problem, the main difficulty was due to the choice of test functions to be used in the associated variational formulation and which are classically some perturbation of the solution to the so-called cellular problem. Indeed, proceeding as usual in homogenization theory, we use energy arguments based on a priori estimates where direct limiting procedure apparently leads to singular behavior. Non local effects appear when these singularities can be overcome, which is usually achieved by using adequate test functions in the variational formulation. Since the fundamental work , an important step was accomplished in this direction in . A slightly different approach uses Dirichlet forms involving non classical measures in the spirit of . However, the main drawback of this method lies in its essential use of the Maximum Principle, which was avoided in for elastic fibers, and later in where the case of spherical symmetry is solved. The asymptotic behavior of highly heterogeneous media has also been considered in the framework of homogenization when the coefficient of one component is vanishing and both components have volumes of unity order: see the derivation of a double porosity model for a single phase flow by and the application of two-scale convergence in order to model diffusion processes in . The paper is organized as follows. Section 2 is devoted to the main notations and to the description of the initial problem. We set the functional framework (16) where the existence and uniqueness of the solution can be established: see and for similar problems. In Section 3, we introduce specific tools to handle the limiting process. This is based on the use of the operators $`G_r`$ defined by (41) which have a localizing effect: this observation motivates the additional assumption (49) on the external sources when the radius of the particles is of critical order $`\epsilon ^3`$ with $`\epsilon `$ denoting the period of the distribution. While passing to the limit, the capacity number $`\gamma _\epsilon `$ defined by (36) appears as the main criterium to describe the limit problem, the relative parts played by the radius of the particles and by the period of the network becoming explicit. Section 4, which is actually the most involving one, deals with the critical case when $`\gamma _\epsilon `$ has a positive and finite limit $`\gamma `$. In this part, where we assume also $`b_\epsilon +\mathrm{}`$ the test functions are a convex combination of the elementary solution of the Laplacian and its transformed by the operator $`G_{r_\epsilon }`$ defined by (41) with $`r=r_\epsilon `$. This choice, which is inspired from , and and has to be compared with , allows to overcome the singular behavior of the energy term when the period $`\epsilon `$ tends to zero. We have to emphasize that this construction highly depends on the geometry of the problem, that is the spherical symmetry . To our knowledge, the generalization to more intricate geometries remains to be done. The resulting model (72)–(75), with the initial value defined after $`u_0`$ in (21) and $`v_0`$ in (23), involves a pairing $`(u,v)`$ which is coupled through a linear operator acting on the difference $`uv`$ by the factor $`4\pi \gamma `$. The case of the infinite capacity, where $`\epsilon ^3<<r_\epsilon <<\epsilon `$, is worked out in Section 5. The proofs are only sketched because the arguments follow the same lines as in Section 4. Let us mention that the singular behavior of the capacity in this case, that is $`\gamma _\epsilon +\mathrm{}`$, forces $`v`$ to coincide with $`u`$. In other words, the infinite capacity prevents the splitting of the distribution, as it did in the critical case. Quite interestingly, the initial value of the global concentration is a convex combination (86) of the initial conditions $`u_0`$ and $`v_0`$; moreover, the mass density of the macroscopic diffusion equation (85) takes both components into account, in accordance with the intuition that the limiting process must lead to a binary mixture. Finally, the case of vanishing capacity is handled in Section 6, that is when $`r_\epsilon <<\epsilon ^3`$. Here, $`v`$ remains constant in time, obviously equal to the initial condition $`v_0`$, while $`u`$ satisfies the diffusion equation (90)–(91) with data independent of the initial condition of the suspension. This can be seen as a proof that when the radius of the particles is too small, then the suspension does not present macroscopic effects, although a corresponding residual concentration, constant in time, should be considered. ## 2 The diffusion problem We consider $`\mathrm{\Omega }𝐑^3`$ a bounded Lipschitz domain occupied by a mixture of two different materials, one of them forming the ambiental connected phase and the other being concentrated in a periodical suspension of small spherical particles. Let us denote $$Y:=(\frac{1}{2},+\frac{1}{2})^3.$$ (1) $$Y_\epsilon ^k:=\epsilon k+\epsilon Y,k𝐙^3.$$ (2) $$𝐙_\epsilon :=\{k𝐙^3,Y_\epsilon ^k\mathrm{\Omega }\},\mathrm{\Omega }_{Y_\epsilon }:=_{k𝐙_\epsilon }Y_\epsilon ^k.$$ (3) The suspension is defined by the following reunion $$D_\epsilon :=_{k𝐙_\epsilon }B(\epsilon k,r_\epsilon ),$$ (4) where $`0<r_\epsilon <<\epsilon `$ and $`B(\epsilon k,r_\epsilon )`$ is the ball of radius $`r_\epsilon `$ centered at $`\epsilon k`$, $`k𝐙_\epsilon `$. Obviously, $$\|D_\epsilon \|0\text{as}\epsilon 0.$$ (5) The fluid domain is given by $$\mathrm{\Omega }_\epsilon =\mathrm{\Omega }D_\epsilon .$$ (6) We also use the following notation for the cylindrical time-domain: $$\mathrm{\Omega }^T:=\mathrm{\Omega }\times ]0,T[;$$ (7) similar definitions for $`\mathrm{\Omega }_\epsilon ^T`$, $`\mathrm{\Omega }_{Y_\epsilon }^T`$ and $`D_\epsilon ^T`$. We consider the problem which governs the diffusion process throughout our binary mixture. Denoting by $`a_\epsilon >0`$ and $`b_\epsilon >0`$ the relative mass density and diffusivity of the suspension, then, assuming without loss of generality that $`\|\mathrm{\Omega }\|=1`$, its non-dimensional form is the following: To find $`u^\epsilon `$ solution of $$\rho ^\epsilon \frac{u^\epsilon }{t}\mathrm{div}(k^\epsilon u^\epsilon )=f^\epsilon \text{in}\mathrm{\Omega }^T$$ (8) $$[u^\epsilon ]_\epsilon =0\text{on}D_\epsilon ^T$$ (9) $$[k^\epsilon u^\epsilon ]_\epsilon n=0\text{on}D_\epsilon ^T$$ (10) $$u^\epsilon =0\text{on}\mathrm{\Omega }^T$$ (11) $$u^\epsilon (0)=u_0^\epsilon \text{in}\mathrm{\Omega }$$ (12) where $`[]_\epsilon `$ is the jump across the interface $`D_\epsilon `$, $`n`$ is the normal on $`D_\epsilon `$ in the outward direction, $`f^\epsilon L^2(0,T;H^1(\mathrm{\Omega }))`$, $`u_0^\epsilon L^2(\mathrm{\Omega })`$ and $$\rho ^\epsilon (x)=\{\begin{array}{ccc}1\hfill & \text{if}\hfill & x\mathrm{\Omega }_\epsilon \hfill \\ a_\epsilon \hfill & \text{if}\hfill & xD_\epsilon \hfill \end{array}$$ (13) $$k^\epsilon (x)=\{\begin{array}{ccc}1\hfill & \text{if}\hfill & x\mathrm{\Omega }_\epsilon \hfill \\ b_\epsilon \hfill & \text{if}\hfill & xD_\epsilon \hfill \end{array}$$ (14) Let $`H_\epsilon `$ be the Hilbert space $`L^2(\mathrm{\Omega })`$ endowed with the scalar product $$(u,v)_{H_\epsilon }:=(\rho ^\epsilon u,v)_\mathrm{\Omega }$$ (15) As $`H_0^1(\mathrm{\Omega })`$ is dense in $`H_\epsilon `$ for any fixed $`\epsilon >0`$, we can set $$H_0^1(\mathrm{\Omega })H_\epsilon H_\epsilon ^{}H^1(\mathrm{\Omega })$$ (16) with continuous embeddings. Now, we can present the variational formulation of the problem (8)-(12). To find $`u^\epsilon L^2(0,T;H_0^1(\mathrm{\Omega }))L^{\mathrm{}}(0,T;H_\epsilon )`$ satisfying (in some sense) the initial condition (12) and the following equation $$\frac{d}{dt}(u^\epsilon ,w)_{H_\epsilon }+(k_\epsilon u^\epsilon ,w)_\mathrm{\Omega }=f^\epsilon ,w\text{in}𝒟^{}(0,T),wH_0^1(\mathrm{\Omega })$$ (17) where $`,`$ denotes the duality product between $`H^1(\mathrm{\Omega })`$ and $`H_0^1(\mathrm{\Omega })`$. ###### Theorem 2.1 Under the above hypotheses and notations, problem (17) has a unique solution. Moreover, $`{\displaystyle \frac{du^\epsilon }{dt}}L^2(0,T;H^1(\mathrm{\Omega }))`$ and hence, $`u^\epsilon `$ is equal almost everywhere to a function of $`C^0([0,T];H_\epsilon )`$; this is the sense of the initial condition (12). In the following we consider that the density of the spherical particles is much higher than that of the surrounding phase. The specific feature of our mixture, which describes the fact that although the volume of the suspension is vanishing its mass is of unity order, is given by: $$\underset{\epsilon 0}{lim}a_\epsilon \|D_\epsilon \|=a>0$$ (18) Regarding the relative diffusivity, we only assume: $$b_\epsilon b_0>0,\epsilon >0.$$ (19) As for the data, we assume that there exist $`fL^2(0,T;H^1(\mathrm{\Omega }))`$ and $`u_0L^2(\mathrm{\Omega })`$ such that $$f^\epsilon f\text{in}L^2(0,T;H^1(\mathrm{\Omega }))$$ (20) $$u_0^\epsilon u_0\text{in}L^2(\mathrm{\Omega })$$ (21) Also, we assume that there exist $`C>0`$ (independent of $`\epsilon `$) and $`v_0L^2(\mathrm{\Omega })`$ for which $$_{D_\epsilon }\|u_0^\epsilon \|^2dxC$$ (22) $$\frac{1}{\|D_\epsilon \|}u_0^\epsilon \chi _{D_\epsilon }v_0\text{in}𝒟^{}(\mathrm{\Omega })$$ (23) where, for any $`D\mathrm{\Omega }`$, we denote $$_Ddx=\frac{1}{\|D\|}_Ddx.$$ ###### Remark 2.2 As $`u_0^\epsilon `$ satisfies (22) then (23) holds at least on some subsequence (see Lemma A-2 ). ###### Proposition 2.3 We have $$u^\epsilon \text{is bounded in }L^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))L^2(0,T;H_0^1(\mathrm{\Omega })).$$ (24) Moreover, there exists $`C>0`$, independent of $`\epsilon `$, such that $$_{D_\epsilon }\|u^\epsilon \|^2dxC\text{a.e. in }[0,T]$$ (25) $$b_\epsilon \|u^\epsilon \|_{L^2(D_\epsilon ^T)}^2C.$$ (26) Proof. Substituting $`w=u^\epsilon `$ in the variational problem (17) and integrating over $`(0,t)`$ for any $`t]0,T[`$, we get: $$\frac{1}{2}\left(\|u^\epsilon (t)\|_{\mathrm{\Omega }_\epsilon }^2+a_\epsilon \|u^\epsilon (t)\|_{D_\epsilon }^2\right)+b_\epsilon _0^t\|u^\epsilon \|_{D_\epsilon }^2𝑑s+_0^t\|u^\epsilon \|_{\mathrm{\Omega }_\epsilon }^2𝑑s=$$ $$=_0^tf^\epsilon (s),u^\epsilon (s)𝑑s+\frac{1}{2}\left(\|u_0^\epsilon \|_{\mathrm{\Omega }_\epsilon }^2+a_\epsilon \|u_0^\epsilon \|_{D_\epsilon }^2\right).$$ Notice that (21) and (22) yield: $$\|u_0^\epsilon \|_{\mathrm{\Omega }_\epsilon }^2+a_\epsilon \|u_0^\epsilon \|_{D_\epsilon }^2\|u_0^\epsilon \|_\mathrm{\Omega }^2+a_\epsilon \|D_\epsilon \|_{D_\epsilon }\|u_0^\epsilon \|^2dxC.$$ Moreover: $$_0^tf^\epsilon (s),u^\epsilon (s)𝑑s_0^t\|f^\epsilon \|_{H^1}\|u^\epsilon \|_\mathrm{\Omega }𝑑s$$ $$_0^t\|f^\epsilon \|_{H^1}\|u^\epsilon \|_{\mathrm{\Omega }_\epsilon }𝑑s+_0^t\|f^\epsilon \|_{H^1}\|u^\epsilon \|_{D_\epsilon }𝑑s$$ $$\frac{1}{2}_0^T\|f^\epsilon \|_{H^1}^2𝑑s+\frac{1}{2}_0^t\|u^\epsilon \|_{\mathrm{\Omega }_\epsilon }^2𝑑s+\frac{1}{2b_\epsilon }_0^T\|f^\epsilon \|_{H^1}^2𝑑s+\frac{b_\epsilon }{2}_0^t\|u^\epsilon \|_{D_\epsilon }^2𝑑s.$$ There results: $$\frac{1}{2}\left(\|u^\epsilon (t)\|_{\mathrm{\Omega }_\epsilon }^2+a_\epsilon \|u^\epsilon (t)\|_{D_\epsilon }^2\right)+\frac{b_\epsilon }{2}_0^t\|u^\epsilon \|_{D_\epsilon }^2𝑑s+\frac{1}{2}_0^t\|u^\epsilon \|_{\mathrm{\Omega }_\epsilon }^2𝑑sC$$ and the proof is completed. ## 3 Specific tools First, we introduce $$_\epsilon =\{R,r_\epsilon <<R<<\epsilon \}$$ that is $`R_\epsilon `$ iff $$\underset{\epsilon 0}{lim}\frac{r_\epsilon }{R}=\underset{\epsilon 0}{lim}\frac{R}{\epsilon }=0.$$ (27) We have to remark that $`_\epsilon `$ is an infinite set, this property being insured by the assumption $`0<r_\epsilon <<\epsilon `$. We denote the domain confined between the spheres of radius $`a`$ and $`b`$ by $$𝒞(a,b):=\{x𝐑^3,a<\|x\|<b\}$$ and correspondingly $$𝒞^k(a,b):=\epsilon k+𝒞(a,b).$$ For any $`R_\epsilon _\epsilon `$, we use the following notations: $$𝒞_\epsilon :=_{k𝐙_\epsilon }𝒞^k(r_\epsilon ,R_\epsilon ),𝒞_\epsilon ^T:=𝒞_\epsilon \times ]0,T[$$ ###### Definition 3.1 For any $`R_\epsilon _\epsilon `$, we define $`w_{R_\epsilon }H_0^1(\mathrm{\Omega })`$ by $`w_{R_\epsilon }(x)`$ $`:=`$ $`\{\begin{array}{c}0\text{in}\mathrm{\Omega }_\epsilon 𝒞_\epsilon ,\\ W_{R_\epsilon }(x\epsilon k)\text{in}𝒞_\epsilon ^k,k𝐙_\epsilon ,\\ 1\text{in}D_\epsilon .\end{array}`$ (31) where $$W_{R_\epsilon }(y)=\frac{r_\epsilon }{(R_\epsilon r_\epsilon )}\left(\frac{R_\epsilon }{\|y\|}1\right)\text{for}y𝒞(r_\epsilon ,R_\epsilon )$$ (32) We have to remark here that $`W_{R_\epsilon }H^1(C(r_\epsilon ,R_\epsilon ))`$ and satisfies the system $`\mathrm{\Delta }W_{R_\epsilon }`$ $`=`$ $`0\text{in}𝒞(r_\epsilon ,R_\epsilon )`$ (33) $`W_{R_\epsilon }`$ $`=`$ $`1\text{for}\|y\|=r_\epsilon `$ (34) $`W_{R_\epsilon }`$ $`=`$ $`0\text{for}\|y\|=R_\epsilon `$ (35) From now on, we denote $$\gamma _\epsilon :=\frac{r_\epsilon }{\epsilon ^3}.$$ (36) ###### Proposition 3.2 For any $`R_\epsilon _\epsilon `$, we have $$\|w_{R_\epsilon }\|_\mathrm{\Omega }C\gamma _\epsilon ^{1/2}$$ (37) $$w_{R_\epsilon }0\text{in}L^2(\mathrm{\Omega }).$$ (38) Proof. First notice that $$\|w_{R_\epsilon }\|_\mathrm{\Omega }=\|w_{R_\epsilon }\|_{𝒞_\epsilon D_\epsilon }\|𝒞_\epsilon D_\epsilon \|^{1/2}C\left(\frac{R_\epsilon }{\epsilon }\right)^{3/2}$$ and $`lim_{\epsilon 0}\frac{R_\epsilon }{\epsilon }=0`$ by assumption (27). As for the rest, direct computation shows $`\|w_{R_\epsilon }\|_\mathrm{\Omega }^2`$ $`=`$ $`{\displaystyle \underset{k𝐙_\epsilon }{}}{\displaystyle _{C_{r_\epsilon ,R_\epsilon }^k}}\|w_{R_\epsilon }\|^2𝑑x`$ $`=`$ $`{\displaystyle \underset{k𝐙_\epsilon }{}}{\displaystyle _0^{2\pi }}𝑑\mathrm{\Phi }{\displaystyle _0^\pi }\mathrm{sin}\mathrm{\Theta }d\mathrm{\Theta }{\displaystyle _{r_\epsilon }^{R_\epsilon }}{\displaystyle \frac{dr}{r^2}}\left({\displaystyle \frac{r_\epsilon R_\epsilon }{R_\epsilon r_\epsilon }}\right)^2`$ $``$ $`C{\displaystyle \frac{\|\mathrm{\Omega }\|}{\epsilon ^3}}\left({\displaystyle \frac{1}{r_\epsilon }}{\displaystyle \frac{1}{R_\epsilon }}\right)\left({\displaystyle \frac{r_\epsilon R_\epsilon }{R_\epsilon r_\epsilon }}\right)^2C{\displaystyle \frac{\gamma _\epsilon }{(1\frac{r_\epsilon }{R_\epsilon })}}`$ and the proof is completed by (27). Lemmas 3.3 and 3.4 below are set without proof since they are a three-dimensional adaptation of Lemmas A.3 and A.4 . ###### Lemma 3.3 For every $`0<r_1<r_2`$ and $`uH^1(C(r_1,r_2))`$, the following estimate holds true: $$\|u\|_{C(r_1,r_2)}^2\frac{4\pi r_1r_2}{r_2r_1}\left\|_{𝐒_{r_2}}ud\sigma _{𝐒_{r_1}}ud\sigma \right\|^2,$$ (39) where $$_{𝐒_r}d\sigma :=\frac{1}{4\pi r^2}_{𝐒_r}d\sigma .$$ ###### Lemma 3.4 There exists a positive constant $`C>0`$ such that: $`(R,\alpha )𝐑^+\times (0,1)`$, $`uH^1(B(0,R))`$, $$_{B(0,R)}\|u_{𝐒_{\alpha R}}ud\sigma \|^2𝑑xC\frac{R^2}{\alpha }\|u\|_{B(0,R)}^2.$$ (40) ###### Definition 3.5 Consider the piecewise constant functions $`G_r:L^2(0,T;H_0^1(\mathrm{\Omega }))L^2(\mathrm{\Omega }^T)`$ defined for any $`r>0`$ by $$G_r(\theta )(x,t)=\underset{k𝐙_\epsilon }{}\left(_{𝐒_r^k}\theta (y,t)d\sigma _y\right)1_{Y_\epsilon ^k}(x)$$ (41) where we denote $$S_r^k=B(\epsilon k,r).$$ (42) ###### Lemma 3.6 If $`R_\epsilon _\epsilon `$, then for every $`\theta L^2(0,T;H_0^1(\mathrm{\Omega }))`$ we have $`\|\theta G_{R_\epsilon }(\theta )\|_{L^2(\mathrm{\Omega }_{Y_\epsilon }^T)}`$ $``$ $`C\left({\displaystyle \frac{\epsilon ^3}{R_\epsilon }}\right)^{1/2}\|\theta \|_{L^2(\mathrm{\Omega }^T)}`$ (43) $`\|\theta G_{r_\epsilon }(\theta )\|_{L^2(D_\epsilon ^T)}`$ $``$ $`Cr_\epsilon \|\theta \|_{L^2(D_\epsilon ^T)}`$ (44) $`\|G_{R_\epsilon }(\theta )G_{r_\epsilon }(\theta )\|_{L^2(\mathrm{\Omega }^T)}`$ $``$ $`C\left({\displaystyle \frac{\epsilon ^3}{r_\epsilon }}\right)^{1/2}\|\theta \|_{L^2(𝒞_\epsilon ^T)}`$ (45) where $`G_{R_\epsilon }(\theta )`$ and $`G_{r_\epsilon }(\theta )`$ are defined following (41). Moreover: $$\|G_{R_\epsilon }(\theta )\|_{L^2(\mathrm{\Omega }^T)}^2=_0^T_{D_\epsilon }\|G_{R_\epsilon }(\theta )\|^2dxdt,\|G_{r_\epsilon }(\theta )\|_{L^2(\mathrm{\Omega }^T)}^2=_0^T_{D_\epsilon }\|G_{r_\epsilon }(\theta )\|^2dxdt.$$ (46) Proof. Notice that by definition: $$\underset{k𝐙_\epsilon }{}_0^T_{Y_\epsilon ^k}\|\theta _{𝐒_{R_\epsilon }^k}\theta d\sigma \|^2𝑑x𝑑t\underset{k𝐙_\epsilon }{}_0^T_{B(\epsilon k,\frac{\epsilon \sqrt{3}}{2})}\|\theta _{𝐒_{R_\epsilon }^k}\theta d\sigma \|^2𝑑x𝑑t$$ where we have used that $$Y_\epsilon ^kB(\epsilon k,\frac{\epsilon \sqrt{3}}{2})$$ for every $`k𝐙_\epsilon `$. We use Lemma 3.4 with $$R=\frac{\epsilon \sqrt{3}}{2},\alpha =\frac{2R_\epsilon }{\epsilon \sqrt{3}}$$ to deduce that $$_{\mathrm{\Omega }_{Y_\epsilon }^T}\|\theta G_{R_\epsilon }(\theta )\|^2𝑑x𝑑tC\left(\frac{\epsilon \sqrt{3}}{2}\right)^2\frac{\epsilon \sqrt{3}}{2R_\epsilon }\underset{k𝐙_\epsilon }{}_0^T_{B(\epsilon k,\frac{\epsilon \sqrt{3}}{2})}\|\theta \|^2𝑑x𝑑t$$ $$C\frac{\epsilon ^3}{R_\epsilon }\underset{k𝐙_\epsilon }{}_0^T_{B(\epsilon k,\frac{\epsilon \sqrt{3}}{2})}\|\theta \|^2𝑑x𝑑tC\frac{\epsilon ^3}{R_\epsilon }_{\mathrm{\Omega }^T}\|\theta \|^2𝑑x𝑑t$$ which shows (43). To establish (44), we recall the definition: $$_{D_\epsilon ^T}\|\theta G_{r_\epsilon }(\theta )\|^2𝑑x𝑑t=\underset{k𝐙_\epsilon }{}_0^T_{B(\epsilon k,r_\epsilon )}\|\theta _{𝐒_{r_\epsilon }^k}\theta d\sigma \|^2𝑑x𝑑t$$ Applying Lemma 3.4 with $`R=r_\epsilon `$ and $`\alpha =1`$, we get the result $`{\displaystyle _{D_\epsilon ^T}}\|\theta G_{r_\epsilon }(\theta )\|^2𝑑x𝑑t`$ $``$ $`Cr_\epsilon ^2{\displaystyle \underset{k𝐙_\epsilon }{}}{\displaystyle _0^T}{\displaystyle _{B(\epsilon k,r_\epsilon )}}\|\theta \|^2𝑑x𝑑tCr_\epsilon ^2{\displaystyle _{D_\epsilon ^T}}\|\theta \|^2𝑑x𝑑t.`$ We come to (45). Indeed, applying Lemma 3.3 and (27): $$_{\mathrm{\Omega }^T}\|G_{R_\epsilon }(\theta )G_{r_\epsilon }(\theta )\|^2𝑑x𝑑t=\underset{k𝐙_\epsilon }{}_0^T_{Y_\epsilon ^k}\|_{𝐒_{R_\epsilon }^k}\theta d\sigma _{𝐒_{r_\epsilon }^k}\theta d\sigma \|^2𝑑y𝑑t$$ $$\underset{k𝐙_\epsilon }{}_{Y_\epsilon ^k}\frac{(R_\epsilon r_\epsilon )}{4\pi R_\epsilon r_\epsilon }𝑑y_0^T_{C_{r_\epsilon ,R_\epsilon }^k}\|\theta \|^2𝑑x𝑑t=\frac{(R_\epsilon r_\epsilon )}{4\pi r_\epsilon R_\epsilon }\underset{k𝐙_\epsilon }{}\epsilon ^3_0^T_{C_{r_\epsilon ,R_\epsilon }^k}\|\theta \|^2𝑑x𝑑t$$ $$=C\epsilon ^3\frac{(R_\epsilon r_\epsilon )}{4\pi r_\epsilon R_\epsilon }_{𝒞_\epsilon ^T}\|\theta \|^2𝑑x𝑑tC\frac{\epsilon ^3}{r_\epsilon }_{𝒞_\epsilon ^T}\|\theta \|^2𝑑x𝑑t.$$ Finally, a direct computation yields (46). ###### Proposition 3.7 If $`R_\epsilon _\epsilon `$, then for any $`\theta L^2(0,T;H_0^1(\mathrm{\Omega }))`$ there holds true: $$_0^T_{D_\epsilon }\|\theta \|^2dxdtC\mathrm{max}(1,\frac{\epsilon ^3}{r_\epsilon })\|\theta \|_{L^2(\mathrm{\Omega }^T)}^2.$$ Proof. We have: $$_0^T_{D_\epsilon }\|\theta \|^2dxdt2_0^T_{D_\epsilon }\|\theta G_{r_\epsilon }(\theta )\|^2dxdt+2_0^T_{D_\epsilon }\|G_{r_\epsilon }(\theta )\|^2dxdt$$ $$=2_0^T_{D_\epsilon }\|\theta G_{r_\epsilon }(\theta )\|^2dxdt+2_{\mathrm{\Omega }^T}\|G_{r_\epsilon }(\theta )\|^2𝑑x𝑑t$$ $$Cr_\epsilon ^2_0^T_{D_\epsilon }\|\theta \|^2dxdt+4_{\mathrm{\Omega }^T}\|G_{r_\epsilon }(\theta )G_{R_\epsilon }(\theta )\|^2𝑑x𝑑t+$$ $$+8_{\mathrm{\Omega }^T}\|G_{R_\epsilon }(\theta )\theta \|^2𝑑x𝑑t+8_{\mathrm{\Omega }^T}\|\theta \|^2𝑑x𝑑t$$ $$Cr_\epsilon ^2_0^T_{D_\epsilon }\|\theta \|^2dxdt+C\frac{\epsilon ^3}{r_\epsilon }_{𝒞_\epsilon ^T}\|\theta \|^2𝑑x𝑑t+$$ $$+C\frac{\epsilon ^3}{R_\epsilon }_{\mathrm{\Omega }^T}\|\theta \|^2𝑑x𝑑t+C_{\mathrm{\Omega }^T}\|\theta \|^2𝑑x𝑑t$$ $$C\left(\frac{\epsilon ^3}{r_\epsilon }+\frac{\epsilon ^3}{R_\epsilon }+1\right)_{\mathrm{\Omega }^T}\|\theta \|^2𝑑x𝑑tC\mathrm{max}(1,\frac{\epsilon ^3}{r_\epsilon })_{\mathrm{\Omega }^T}\|\theta \|^2𝑑x𝑑t$$ ###### Remark 3.8 Using the Mean Value Theorem, we easily find that $$\|G_{r_\epsilon }(\phi )\phi \|_{L^{\mathrm{}}(𝒞_\epsilon D_\epsilon )}2R_\epsilon \|\phi \|_{L^{\mathrm{}}(\mathrm{\Omega })},\phi 𝒟(\mathrm{\Omega }).$$ ###### Definition 3.9 Let $`M_{D_\epsilon }:L^2(0,T;C_c(\mathrm{\Omega }))L^2(\mathrm{\Omega }^T)`$ be defined by $$M_{D_\epsilon }(\phi )(x,t):=\underset{k𝐙_\epsilon }{}\left(_{Y_\epsilon ^k}\phi (y,t)dy\right)\mathrm{\hspace{0.33em}1}_{B(\epsilon k,r_\epsilon )}(x).$$ ###### Lemma 3.10 For any $`\phi L^2(0,T;C_c(\mathrm{\Omega }))`$, we have: $$\underset{\epsilon 0}{lim}_0^T_{D_\epsilon }\|\phi M_{D_\epsilon }(\phi )\|^2dxdt=0.$$ Proof. Notice that $$_{D_\epsilon }\|\phi M_{D_\epsilon }(\phi )\|^2dx=\frac{1}{\|D_\epsilon \|}\underset{k𝐙_\epsilon }{}_{B(\epsilon k,r_\epsilon )}\|\phi _{Y_\epsilon ^k}\phi dy\|^2𝑑x.$$ As $`\mathrm{card}(𝐙_\epsilon ){\displaystyle \frac{\|\mathrm{\Omega }\|}{\epsilon ^3}},\text{then}\|B(0,r_\epsilon )\|{\displaystyle \frac{\mathrm{card}(𝐙_\epsilon )}{\|D_\epsilon \|}}\|\mathrm{\Omega }\|=1`$ and by the uniform continuity of $`\phi `$ on $`\mathrm{\Omega }`$ it follows the convergence to $`0`$ a.e. on $`[0,T]`$. Lebesgue’s dominated convergence theorem achieves the result. ## 4 Homogenization of the case $`r_\epsilon =𝒪(\epsilon ^3)`$ The present critical radius case is described by $$\underset{\epsilon 0}{lim}\gamma _\epsilon =\gamma ]0,+\mathrm{}[.$$ (47) Its homogenization process is the most involving one. That is why we start the homogenization study of our problem with this case, under the condition $$\underset{\epsilon 0}{lim}b_\epsilon =+\mathrm{}$$ (48) We also assume that $`f^\epsilon `$ has the following additional property: $$\{\begin{array}{c}R_\epsilon _\epsilon \text{ and}gL^2(0,T;H^1(\mathrm{\Omega }))\text{for which}\hfill \\ \\ f^\epsilon ,w_{R_\epsilon }\phi g,\phi \text{in}𝒟^{}(0,T),\phi 𝒟(\mathrm{\Omega })\hfill \end{array}$$ (49) (see for a certain type of functions $`f^\epsilon `$ which satisfy (49)). ###### Remark 4.1 Notice that due to (47), in this case Proposition 3.7 reads $$\phi L^2(0,T;H_0^1(\mathrm{\Omega })),_0^T_{D_\epsilon }\|\phi \|^2dxdtC\|\phi \|_{L^2(\mathrm{\Omega }^T)}^2.$$ (50) A preliminary result is the following: ###### Proposition 4.2 There exist $`uL^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))L^2(0,T;H_0^1(\mathrm{\Omega }))`$ and $`vL^2(\mathrm{\Omega }^T)`$ such that, on some subsequence, $$u^\epsilon \stackrel{}{}u\text{in}L^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))$$ (51) $$u^\epsilon u\text{in}L^2(0,T;H_0^1(\mathrm{\Omega }))$$ (52) $$G_{R_\epsilon }(u^\epsilon )u\text{in}L^2(\mathrm{\Omega }^T)$$ (53) $$G_{r_\epsilon }(u^\epsilon )v\text{in}L^2(\mathrm{\Omega }^T)$$ (54) Moreover, we have $$\underset{\epsilon 0}{lim}_0^T_{D_\epsilon }\|u^\epsilon G_{r_\epsilon }(u^\epsilon )\|^2dxdt=0$$ (55) Proof. From (24), we get, on some subsequence, the convergences (51) and (52). Moreover, we have: $$\|uG_{R_\epsilon }(u^\epsilon )\|_{\mathrm{\Omega }^T}^2=\|u\|_{\mathrm{\Omega }^T\mathrm{\Omega }_{Y_\epsilon }^T}^2+\|uG_{R_\epsilon }(u^\epsilon )\|_{\mathrm{\Omega }_{Y_\epsilon }^T}^2$$ (56) where: $$\|uG_{R_\epsilon }(u^\epsilon )\|_{\mathrm{\Omega }_{Y_\epsilon }^T}\|uu^\epsilon \|_{\mathrm{\Omega }_{Y_\epsilon }^T}+\|u^\epsilon G_{R_\epsilon }(u^\epsilon )\|_{\mathrm{\Omega }_{Y_\epsilon }^T}$$ (57) $$\|uu^\epsilon \|_{\mathrm{\Omega }^T}+\|u^\epsilon G_{R_\epsilon }(u^\epsilon )\|_{\mathrm{\Omega }_{Y_\epsilon }^T}$$ and (43) yields: $$\|u^\epsilon G_{R_\epsilon }(u^\epsilon )\|_{\mathrm{\Omega }_{Y_\epsilon }^T}^2C\frac{\epsilon ^3}{R_\epsilon }\|u^\epsilon \|_{\mathrm{\Omega }^T}^2=C\frac{\epsilon ^3}{r_\epsilon }\frac{r_\epsilon }{R_\epsilon }\|u^\epsilon \|_{\mathrm{\Omega }^T}^2C\frac{r_\epsilon }{R_\epsilon }$$ and thus: $$\underset{\epsilon 0}{lim}\|u^\epsilon G_{R_\epsilon }(u^\epsilon )\|_{\mathrm{\Omega }_{Y_\epsilon }^T}^2=0.$$ As (52) implies that $$u^\epsilon u\text{in}L^2(\mathrm{\Omega }^T)$$ (58) the right-hand side of (57) tends to zero as $`\epsilon 0`$, that is: $$\underset{\epsilon 0}{lim}\|uG_{R_\epsilon }(u^\epsilon )\|_{\mathrm{\Omega }_{Y_\epsilon }^T}=0.$$ After substitution into the right-hand side of (56), and taking into account that $$\underset{\epsilon 0}{lim}\|\mathrm{\Omega }^T\mathrm{\Omega }_{Y_\epsilon }^T\|=0,$$ we obtain (53), that is, $$G_{R_\epsilon }(u^\epsilon )u\text{in}L^2(\mathrm{\Omega }^T).$$ (59) In order to prove (54), we see that $$\begin{array}{c}\|G_{r_\epsilon }(u^\epsilon )\|_{L^2(\mathrm{\Omega }^T)}\|G_{r_\epsilon }(u^\epsilon )G_{R_\epsilon }(u^\epsilon )\|_{L^2(\mathrm{\Omega }^T)}+\|G_{R_\epsilon }(u^\epsilon )\|_{L^2(\mathrm{\Omega }^T)}\hfill \\ \\ \frac{^C}{\gamma _\epsilon ^{1/2}}\|u^\epsilon \|_{L^2(\mathrm{\Omega }^T)}+CC.\hfill \end{array}$$ (60) Moreover, recall that from (44) we have, taking into account (26): $$_0^T_{D_\epsilon }\|u^\epsilon G_{r_\epsilon }(u^\epsilon )\|^2dxdtCr_\epsilon ^2_0^T_{D_\epsilon }\|u^\epsilon \|^2dxdt\frac{C}{\gamma _\epsilon b_\epsilon }0$$ (61) and the proof is completed. ###### Proposition 4.3 For any $`\phi L^2(0,T;C_c(\mathrm{\Omega }))`$, we have: $$\underset{\epsilon 0}{lim}_0^T_{D_\epsilon }u^\epsilon \phi dxdt=_{\mathrm{\Omega }^T}v\phi 𝑑x𝑑t.$$ (62) Proof. We have: $$\begin{array}{c}_0^T_{D_\epsilon }u^\epsilon \phi dxdt=_0^T_{D_\epsilon }(u^\epsilon G_{r_\epsilon }(u^\epsilon ))\phi dxdt+\hfill \\ \\ +_0^T_{D_\epsilon }G_{r_\epsilon }(u^\epsilon )(\phi M_{D_\epsilon }(\phi ))dxdt+_0^T_{D_\epsilon }G_{r_\epsilon }(u^\epsilon )M_{D_\epsilon }(\phi )dxdt\hfill \end{array}$$ (63) The first right-hand term tends to zero thanks to (55) in Proposition 4.2. The second one tends also to zero thanks to Lemma 3.10. The last term is handled as follows: $$_0^T_{D_\epsilon }G_{r_\epsilon }(u^\epsilon )M_{D_\epsilon }(\phi )dxdt=\lambda _\epsilon \underset{k𝐙_\epsilon }{}_0^T_{Y_\epsilon ^k}\left(_{𝐒_\epsilon ^k}u^\epsilon d\sigma \right)\phi 𝑑x𝑑t=\lambda _\epsilon _{\mathrm{\Omega }^T}\phi G_{r_\epsilon }(u^\epsilon )𝑑x𝑑t$$ where $$\lambda _\epsilon :=\frac{\|B(0,r_\epsilon )\|}{\epsilon ^3\|D_\epsilon \|}1\text{as}\|\mathrm{\Omega }\|=1.$$ The proof is completed by (54). ###### Proposition 4.4 For any $`\phi L^2(0,T;H_0^1(\mathrm{\Omega }))`$, we have $$_0^T_{D_\epsilon }u^\epsilon \phi dxdt_{\mathrm{\Omega }^T}v\phi 𝑑x𝑑t.$$ (64) Proof. In the light of proposition 4.3, we have to prove that the left-hand side term is continuous in the corresponding norm. This can be obtained as follows: $$\left\|_0^T_{D_\epsilon }u^\epsilon \phi dxdt\right\|\left(_0^T_{D_\epsilon }\|u^\epsilon \|^2dxdt\right)^{1/2}\left(_0^T_{D_\epsilon }\|\phi \|^2dxdt\right)^{1/2}$$ $$C\|\phi \|_{L^2(0,T;H_0^1(\mathrm{\Omega }))}^2,$$ where we used (25) and (50). ###### Proposition 4.5 Let for any $`R_\epsilon _\epsilon `$ and $`\phi ,\psi 𝒟(\mathrm{\Omega })`$ $$\mathrm{\Phi }^\epsilon =(1w_{R_\epsilon })\phi +w_{R_\epsilon }G_{r_\epsilon }(\psi )$$ (65) Then, for any $`\eta 𝒟([0,T[)`$, we have $$\underset{\epsilon 0}{lim}\|\mathrm{\Phi }^\epsilon \phi \|_\mathrm{\Omega }=0$$ (66) $$\underset{\epsilon 0}{lim}_{\mathrm{\Omega }^T}\rho ^\epsilon u^\epsilon \mathrm{\Phi }^\epsilon \eta ^{}(t)𝑑x𝑑t=_{\mathrm{\Omega }^T}u\phi \eta ^{}(t)𝑑x𝑑t+a_{\mathrm{\Omega }^T}v\psi \eta ^{}(t)𝑑x𝑑t.$$ (67) Proof. The property (66) is an immediate consequence of (38) and of the uniform boundness of $`G_{r_\epsilon }(\psi )`$ in $`L^{\mathrm{}}(\mathrm{\Omega })`$. For the second property, let us notice that $$_{\mathrm{\Omega }^T}\rho ^\epsilon u^\epsilon \mathrm{\Phi }^\epsilon \eta ^{}(t)𝑑x𝑑t=_0^T_\mathrm{\Omega }\chi _{\mathrm{\Omega }_\epsilon }u^\epsilon \mathrm{\Phi }^\epsilon (x)\eta ^{}(t)𝑑x𝑑t$$ $$+a_\epsilon _0^T_{D_\epsilon }u^\epsilon G_{r_\epsilon }(\psi )\eta ^{}(t)𝑑x𝑑t.$$ As we obviously have $$\underset{\epsilon 0}{lim}_0^T_\mathrm{\Omega }\chi _{\mathrm{\Omega }_\epsilon }u^\epsilon \mathrm{\Phi }^\epsilon (x)\eta ^{}(t)𝑑x𝑑t=_{\mathrm{\Omega }^T}u\phi \eta ^{}(t)𝑑x𝑑t,$$ it remains to study $$a_\epsilon _0^T_{D_\epsilon }u^\epsilon G_{r_\epsilon }(\psi )\eta ^{}(t)𝑑x𝑑t=a_\epsilon \|D_\epsilon \|_0^T_{D_\epsilon }u^\epsilon G_{r_\epsilon }(\psi )\eta ^{}(t)dxdt.$$ Using (62) and the uniform continuity of $`\psi `$, we get $$\underset{\epsilon 0}{lim}a_\epsilon _0^T_{D_\epsilon }u^\epsilon G_{r_\epsilon }(\psi )\eta ^{}(t)𝑑x𝑑t=a_{\mathrm{\Omega }^T}v\psi \eta ^{}𝑑x𝑑t.$$ ###### Proposition 4.6 If $`\mathrm{\Phi }^\epsilon `$ is defined like in Proposition 4.5, then we have $$\underset{\epsilon 0}{lim}_0^T_\mathrm{\Omega }u^\epsilon \mathrm{\Phi }^\epsilon \eta (t)𝑑x𝑑t=_{\mathrm{\Omega }^T}u\phi \eta (t)𝑑x𝑑t+4\pi \gamma _{\mathrm{\Omega }^T}(vu)(\psi \phi )\eta (t)𝑑x𝑑t$$ (68) Proof. First consider $$_0^T_{\mathrm{\Omega }_\epsilon }u^\epsilon \mathrm{\Phi }^\epsilon dxdt$$ which reduces to $$_0^T_{\mathrm{\Omega }_\epsilon 𝒞_\epsilon }u^\epsilon \phi \eta dxdt+_0^T_{𝒞_\epsilon }u^\epsilon \mathrm{\Phi }^\epsilon \eta dxdt.$$ Lebesgue’s dominated convergence theorem yields $`\phi 1_{\mathrm{\Omega }_\epsilon 𝒞_\epsilon }\phi `$ in $`L^2(\mathrm{\Omega })`$. Thus, taking (52) into account $$_0^T_\mathrm{\Omega }u^\epsilon \phi \eta \chi _{\mathrm{\Omega }_\epsilon 𝒞_\epsilon }dxdt_{\mathrm{\Omega }^T}u\phi \eta dxdt.$$ Now, we come to the remaining part, namely $$\begin{array}{ccc}& & _0^T_{𝒞\epsilon }u^\epsilon \mathrm{\Phi }^\epsilon \eta (t)𝑑x𝑑t=_0^T_{𝒞\epsilon }(1w_{R_\epsilon })u^\epsilon \phi \eta dxdt\hfill \\ & & +_0^T_{𝒞\epsilon }u^\epsilon w_{R_\epsilon }(G_{r_\epsilon }(\psi )\phi )𝑑x𝑑t\hfill \\ & :=& I_1+I_2\hfill \end{array}$$ (69) In the first integral, as $`\chi _{𝒞_\epsilon }\phi 0`$ in $`L^2(\mathrm{\Omega }^T)`$, $`u^\epsilon u`$ in $`L^2(\mathrm{\Omega }^T)`$ and $`(1w_{R_\epsilon })`$ is obviously bounded, we easily find that $`I_1`$ tends to zero. In order to study $`I_2`$, let us notice that $$\begin{array}{ccc}I_2\hfill & =& _0^T_{𝒞\epsilon }u^\epsilon w_{R_\epsilon }(G_{r_\epsilon }(\phi )\phi )\eta 𝑑x𝑑t+\hfill \\ & & +_0^T_{𝒞\epsilon }u^\epsilon w_{R_\epsilon }(G_{r_\epsilon }(\psi )G_{r_\epsilon }(\phi ))\eta 𝑑x𝑑t\hfill \end{array}$$ (70) The first term in the right-hand side of (70) may be estimated by $$\|_0^T_{𝒞\epsilon }u^\epsilon w_{R_\epsilon }(\phi G_{r_\epsilon }(\phi ))\eta 𝑑x𝑑t\|\|u^\epsilon \|_{\mathrm{\Omega }^T}\|w_{R_\epsilon }\eta \|_{\mathrm{\Omega }^T}\|\phi G_{r_\epsilon }(\phi )\|_{L^{\mathrm{}}(𝒞_\epsilon )}.$$ (71) As $`(w_{R_\epsilon })`$ is bounded in $`H^1(\mathrm{\Omega })`$ (see Proposition 3.2), the right hand side of (71) tends to zero by Remark 3.8. Going back to the second term in the right hand side of (70), we may write $`{\displaystyle _0^T}{\displaystyle _{𝒞\epsilon }}u^\epsilon w_{R_\epsilon }(G_{r_\epsilon }(\psi )G_{r_\epsilon }(\phi ))\eta (t)𝑑x𝑑t`$ $`=`$ $`{\displaystyle \underset{k𝐙_\epsilon }{}}({\displaystyle }_{𝐒_{r_\epsilon }^k}\psi d\sigma {\displaystyle }_{𝐒_{r_\epsilon }^k}\phi d\sigma ){\displaystyle _0^{2\pi }}𝑑\mathrm{\Phi }{\displaystyle _0^\pi }\mathrm{sin}\mathrm{\Theta }d\mathrm{\Theta }{\displaystyle _{r_\epsilon }^{R_\epsilon }}({\displaystyle _0^T}{\displaystyle \frac{u^\epsilon }{r}}\|_{𝒞^k(r_\epsilon ,R_\epsilon )}\eta (t)dt){\displaystyle \frac{dW_{R_\epsilon }}{dr}}r^2𝑑r`$ $`=`$ $`{\displaystyle \frac{r_\epsilon R_\epsilon }{(R_\epsilon r_\epsilon )}}{\displaystyle \underset{k𝐙_\epsilon }{}}\left({\displaystyle }_{𝐒_{r_\epsilon }^k}\psi d\sigma {\displaystyle }_{𝐒_{r_\epsilon }^k}\phi d\sigma \right){\displaystyle _{𝐒_1}}{\displaystyle _0^T}(u^\epsilon \|_{\|x\epsilon k\|=r_\epsilon }u^\epsilon \|_{\|x\epsilon k\|=R_\epsilon })\eta (t)𝑑t𝑑\sigma _1`$ $`=`$ $`{\displaystyle \frac{4\pi r_\epsilon R_\epsilon }{\epsilon ^3(R_\epsilon r_\epsilon )}}{\displaystyle _{\mathrm{\Omega }^T}}(G_{r_\epsilon }(u^\epsilon )G_{R_\epsilon }(u^\epsilon ))(G_{r_\epsilon }(\psi )G_{r_\epsilon }(\phi ))\eta (t)𝑑x𝑑t`$ from which we infer that $`I_2`$ is converging to $$4\pi \gamma _{\mathrm{\Omega }^T}(vu)(\psi \phi )\eta (t)𝑑x𝑑t$$ and the proof is completed. We are in the position to state our main result: ###### Theorem 4.7 The limits $`uL^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))L^2(0,T;H_0^1(\mathrm{\Omega }))`$ and $`vL^2(\mathrm{\Omega }^T)`$ of (51)–(54) verify (in a weak sense) the following problem: $`{\displaystyle \frac{u}{t}}\mathrm{\Delta }u+4\pi \gamma (uv)`$ $`=`$ $`(fg)\text{in}\mathrm{\Omega }^T,`$ (72) $`a{\displaystyle \frac{v}{t}}+4\pi \gamma (vu)`$ $`=`$ $`g\text{in}\mathrm{\Omega }^T,`$ (73) $`u(0)`$ $`=`$ $`u_0\text{in}\mathrm{\Omega }`$ (74) $`v(0)`$ $`=`$ $`v_0\text{in}\mathrm{\Omega }`$ (75) Moreover, there holds $`uC^0([0,T];L^2(\mathrm{\Omega }))`$ and $`vC^0([0,T];H^1(\mathrm{\Omega }))`$; these are the senses of (74) and (75). ###### Remark 4.8 As the problem (72)–(75) has a unique weak solution, the convergences in Proposition 4.2 hold on the whole sequence. Proof of Theorem 4.7. We set in (17) $`w=\mathrm{\Phi }^\epsilon `$ where $`\mathrm{\Phi }^\epsilon `$ is defined like in lemma 4.5. Then, by multiplying (17) by $`\eta 𝒟([0,T[)`$ and integrating it over $`[0,T]`$ we get $$_{\mathrm{\Omega }^T}\rho ^\epsilon u^\epsilon \mathrm{\Phi }^\epsilon \eta ^{}𝑑x𝑑t+_{\mathrm{\Omega }^T}k^\epsilon u^\epsilon (\mathrm{\Phi }^\epsilon )\eta 𝑑x𝑑t=_0^Tf^\epsilon ,\mathrm{\Phi }^\epsilon \eta 𝑑t+_\mathrm{\Omega }\rho ^\epsilon u_0^\epsilon \mathrm{\Phi }^\epsilon \eta (0)𝑑x.$$ (76) Then, the left-hand side tends to $$\begin{array}{c}_{\mathrm{\Omega }^T}u\phi \eta ^{}𝑑x𝑑ta_{\mathrm{\Omega }^T}v\phi \eta ^{}𝑑x𝑑t+_{\mathrm{\Omega }^T}u\phi \eta dxdt+\\ \\ +4\pi \gamma _{\mathrm{\Omega }^T}(vu)(\psi \phi )\eta 𝑑x𝑑t.\end{array}$$ (77) This is a direct consequence of Proposition 4.6 together with the remark that $$_0^T_{D_\epsilon }u^\epsilon \mathrm{\Phi }^\epsilon dxdt=0$$ since $`\mathrm{\Phi }^\epsilon `$ is constant on every $`B(\epsilon k,r_\epsilon )`$, $`k𝐙_\epsilon `$. As for the right-hand side, we have $$_0^Tf^\epsilon ,\mathrm{\Phi }^\epsilon \eta 𝑑t=_0^Tf^\epsilon ,(1w_{R_\epsilon })\phi \eta 𝑑t+_0^Tf^\epsilon ,w_{R_\epsilon }G_{r_\epsilon }(\psi )\eta 𝑑t$$ and, with hypothesis (49), $$_0^Tf^\epsilon ,(1w_{R_\epsilon })\phi \eta 𝑑t_0^Tfg,\phi \eta 𝑑t.$$ Moreover, $$_0^Tf^\epsilon ,w_{R_\epsilon }G_{r_\epsilon }(\psi )\eta 𝑑t=_0^Tf^\epsilon ,w_{R_\epsilon }(G_{r_\epsilon }(\psi )\psi )\eta 𝑑t+_0^Tf^\epsilon ,w_{R_\epsilon }\psi \eta 𝑑t$$ with $$\left\|_0^Tf^\epsilon ,w_{R_\epsilon }(G_{r_\epsilon }(\psi )\psi )\eta 𝑑t\right\|_0^T\|f^\epsilon \|_{H^1}\|w_{R_\epsilon }(G_{r_\epsilon }(\psi )\psi )\|_{H_0^1(\mathrm{\Omega })}.$$ As we have $$\|w_{R_\epsilon }(G_{r_\epsilon }(\psi )\psi )\|_{H_0^1(\mathrm{\Omega })}=\|(w_{R_\epsilon }(G_{r_\epsilon }(\psi )\psi ))\|_\mathrm{\Omega }$$ $$\|w_{R_\epsilon }\|_{𝒞_\epsilon }\|G_{r_\epsilon }(\psi )\psi \|_{L^{\mathrm{}}(𝒞_\epsilon )}+\|\psi \|_{𝒞_\epsilon D_\epsilon }$$ Remark 3.8 and (37) obviously yield $$\underset{\epsilon 0}{lim}\|w_{R_\epsilon }(G_{r_\epsilon }(\psi )\psi )\|_{H_0^1(\mathrm{\Omega })}=0.$$ The assumption (20) on $`f^\epsilon `$ implies that $`\|f^\epsilon \|_{H^1}C`$ and thus $$\underset{\epsilon 0}{lim}_0^Tf^\epsilon ,w_{R_\epsilon }(G_{r_\epsilon }(\psi )\psi )\eta 𝑑t=0.$$ We conclude thanks to hypothesis (49) that $$\underset{\epsilon 0}{lim}_0^Tf^\epsilon ,w_{R_\epsilon }\psi \eta 𝑑t=_0^Tg,\psi \eta 𝑑t.$$ Finally: $$\underset{\epsilon 0}{lim}_0^Tf^\epsilon ,\mathrm{\Phi }^\epsilon \eta 𝑑t=_0^Tfg,\phi \eta 𝑑t+_0^Tg,\psi \eta 𝑑t.$$ We get $$_\mathrm{\Omega }\rho ^\epsilon u_0^\epsilon \mathrm{\Phi }^\epsilon \eta (0)𝑑x=_{\mathrm{\Omega }_\epsilon }u_0^\epsilon \mathrm{\Phi }^\epsilon \eta (0)𝑑x+a_\epsilon _{D_\epsilon }u_0^\epsilon G_{r_\epsilon }(\psi )\eta (0)𝑑x.$$ Using the hypotheses (21)–(23) on $`u_0^\epsilon `$, we pass to the limit and with the same arguments as above we obtain $$\underset{\epsilon 0}{lim}_\mathrm{\Omega }\rho ^\epsilon u_0^\epsilon \mathrm{\Phi }^\epsilon \eta (0)𝑑x=\eta (0)_\mathrm{\Omega }(u_0\phi +av_0\psi )𝑑x$$ which achieves the proof. ## 5 Homogenization in the case $`\epsilon ^3<<r_\epsilon <<\epsilon `$ In this section, we fix some $`R_\epsilon _\epsilon `$. ###### Remark 5.1 Notice that in this case Proposition 3.7 also reads $$_0^T_{D_\epsilon }\|\phi \|^2dxdtC\|\phi \|_{L^2(\mathrm{\Omega }^T)}^2,\phi L^2(0,T;H_0^1(\mathrm{\Omega })).$$ (78) In the present case, Proposition 2.3 and Lemma 3.6 imply in a straightforward manner the result corresponding to Proposition 4.2. ###### Proposition 5.2 There exists $`uL^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))L^2(0,T;H_0^1(\mathrm{\Omega }))`$ such that, on some subsequence, $$u^\epsilon \stackrel{}{}u\text{in}L^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))$$ (79) $$u^\epsilon u\text{in}L^2(0,T;H_0^1(\mathrm{\Omega }))$$ (80) $$G_{R_\epsilon }(u^\epsilon )u\text{in}L^2(\mathrm{\Omega }^T)$$ (81) $$G_{r_\epsilon }(u^\epsilon )u\text{in}L^2(\mathrm{\Omega }^T)$$ (82) Moreover, we have $$\underset{\epsilon 0}{lim}_0^T_{D_\epsilon }\|u^\epsilon G_{r_\epsilon }(u^\epsilon )\|^2dxdt=0$$ (83) In the light of Remark 5.1, we prove as in the previous section: ###### Proposition 5.3 For any $`\phi L^2(0,T;H_0^1(\mathrm{\Omega }))`$, we have $$_0^T_{D_\epsilon }u^\epsilon \phi dxdt_{\mathrm{\Omega }^T}u\phi 𝑑x𝑑t.$$ (84) The homogenization result obtained in this case follows. ###### Theorem 5.4 The limit $`uL^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))L^2(0,T;H_0^1(\mathrm{\Omega }))`$ of (79)–(82) is the only solution of $`(1+a){\displaystyle \frac{u}{t}}\mathrm{\Delta }u=f\text{in}\mathrm{\Omega }^T,`$ (85) $`u(0)={\displaystyle \frac{1}{(1+a)}}u_0+{\displaystyle \frac{a}{(1+a)}}v_0\text{in}\mathrm{\Omega }`$ (86) Moreover, the convergences in Proposition 5.2 hold on the whole sequence and $`uC^0([0,T];L^2(\mathrm{\Omega }))`$, this being the sense of (86). Proof. The proof of (85) is similar to the corresponding one of the Theorem 4.7. The test function $`\mathrm{\Phi }^\epsilon `$ is given by $$\mathrm{\Phi }^\epsilon =(1w_{R_\epsilon })\phi +w_{R_\epsilon }G_{r_\epsilon }(\phi ),\phi 𝒟(\mathrm{\Omega }).$$ The only interesting convergences are the following two: $$\left\|_{𝒞_\epsilon ^T}u^\epsilon (w_{R_\epsilon })(G_{r_\epsilon }(\phi )\phi )𝑑x𝑑t\right\|C\|u^\epsilon \|_{\mathrm{\Omega }^T}\|w_{R_\epsilon }\|_{\mathrm{\Omega }^T}\left\|G_{r_\epsilon }(\phi )\phi \right\|_{L^{\mathrm{}}(𝒞_\epsilon ^T)}$$ $$C\gamma _\epsilon ^{1/2}R_\epsilon =C\left(\frac{r_\epsilon }{\epsilon }\right)^{1/2}\left(\frac{R_\epsilon }{\epsilon }\right)0$$ $$\left\|_0^Tf^\epsilon ,w_{R_\epsilon }(G_{r_\epsilon }(\phi )\phi )\right\|C\left\|(G_{r_\epsilon }(\phi )\phi )w_{R_\epsilon }\right\|_{L^2(𝒞_\epsilon ^T)}+C\left\|w_{R_\epsilon }\phi \right\|_{L^2(𝒞_\epsilon ^TD_\epsilon ^T)}$$ $$C\left\|\phi \right\|_{L^{\mathrm{}}(\mathrm{\Omega })}\left(\gamma _\epsilon ^{1/2}R_\epsilon +\|𝒞_\epsilon D_\epsilon \|^{1/2}\right)0,$$ where we have used the a priori estimates of Proposition 2.3, Remark 3.8 and Proposition 3.2. Using Proposition 5.2 and hypotheses (21)–(23), we obtain with the same argument as before $$\underset{\epsilon 0}{lim}_\mathrm{\Omega }\rho ^\epsilon u_0^\epsilon \mathrm{\Phi }^\epsilon \eta (0)𝑑x=\eta (0)_\mathrm{\Omega }(u_0+av_0)\phi 𝑑x$$ which achieves the proof. ## 6 Homogenization in the case $`r_\epsilon <<\epsilon ^3`$. As in this case $`\gamma _\epsilon 0`$, we only can prove: ###### Theorem 6.1 There exists $`uL^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))L^2(0,T;H_0^1(\mathrm{\Omega }))`$ such that $`u^\epsilon \stackrel{}{}u`$ in $`L^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))`$ (87) $`u^\epsilon u`$ in $`L^2(0,T;H_0^1(\mathrm{\Omega }))`$ (88) $`{\displaystyle \frac{1}{\|D_\epsilon \|}}u^\epsilon \chi _{D_\epsilon }v_0`$ in $`𝒟^{}(\mathrm{\Omega })\text{a.e.}t[0,T]`$ (89) where $`u`$ is the only solution of the following problem: $`{\displaystyle \frac{u}{t}}\mathrm{\Delta }u`$ $`=`$ $`f\text{in}\mathrm{\Omega }^T`$ (90) $`u(0)`$ $`=`$ $`u_0\text{in}\mathrm{\Omega }`$ (91) Proof. The convergences (87)–(88) hold on some subsequences; they are insured by Proposition 2.3. We have to remark that (25) is the hypothesis which insures the existence of $`vL^{\mathrm{}}(0,T;L^2(\mathrm{\Omega }))`$ which satisfies $$\frac{1}{\|D_\epsilon \|}u^\epsilon \chi _{D_\epsilon }v\text{in}𝒟^{}(\mathrm{\Omega })\text{a.e.}t[0,T]$$ (see Lemma A-2 ); we have to prove that $`v=v_0`$. Acting as usual, we take $$\mathrm{\Phi }^\epsilon =(1w_{R_\epsilon })\phi +w_{R_\epsilon }G_{r_\epsilon }(\psi )$$ (92) for some $`R_\epsilon _\epsilon `$ and $`\phi ,\psi 𝒟(\mathrm{\Omega })`$. Notice that in this case we have $$\mathrm{\Phi }^\epsilon \phi \text{in}H_0^1(\mathrm{\Omega })$$ (93) because obviously $`w_{R_\epsilon }0`$ in $`H_0^1(\mathrm{\Omega })`$. Passing to the limit in the variational formulation, we obtain in a straightforward manner $$_{\mathrm{\Omega }^T}u\phi \eta ^{}𝑑x𝑑ta_{\mathrm{\Omega }^T}v\psi \eta ^{}𝑑x𝑑t+_{\mathrm{\Omega }^T}u\phi \eta dxdt=_0^Tf,\phi \eta 𝑑t+$$ $$+(_\mathrm{\Omega }u_0\phi dx+a_\mathrm{\Omega }v_0\psi dx)\eta (0),\eta 𝒟([0,T[)$$ Setting $`\phi =0`$, we find that $`v`$ is independent of $`t`$ and that $`vC^0([0,T];L^2(\mathrm{\Omega }))`$, which achieves $`v=v_0`$. Then, setting $`\psi =0`$, we prove (90) and (91), the last one holding also in the sense of $`C^0([0,T];L^2(\mathrm{\Omega }))`$. Acknowledgements. This work was done during the visit of F. Bentalha and D. Polişevschi at the I.R.M.A.R.’s Department of Mechanics (University of Rennes 1) whose support is gratefully acknowledged. Also, this work corresponds to a part of the C.N.C.S.I.S. Research Program 33079-2004. * University of Batna, Department of Mathematics, Batna, Algeria, \ Université de Rennes1, I.R.M.A.R, Campus de Beaulieu, 35042 Rennes Cedex (France) \* I.M.A.R., P.O. Box 1-764, Bucharest (Romania).
warning/0506/cond-mat0506646.html	ar5iv	text	# Direct Observation of Condon Domains in Silver by Hall Probes ## Abstract Using a set of micro Hall probes for the detection of the local induction, the inhomogeneous Condon domain structure has been directly observed at the surface of a pure silver single crystal under strong Landau quantization in magnetic fields up to 10 T. The inhomogeneous induction occurs in the theoretically predicted part of the $`HT`$ Condon domain phase diagram. Information about size, shape and orientation of the domains is obtained by analyzing Hall probes placed along and across the long sample axis and by tilting the sample. On a beryllium surface the induction inhomogeneity is almost absent although the expected induction splitting here is at least ten times higher than in silver. Periodic formation and disappearance of a phase with diamagnetic and paramagnetic domains was predicted by Condon CONDON to occur in a normally nonmagnetic metal in strong magnetic fields ($`H`$) at low temperatures ($`T`$). The domains arise under the condition $`\chi =\mu _0M/B>1`$, where $`M`$ is the oscillating magnetization of the electrons due to Landau quantization and $`B=\mu _0(H+M)`$ the total induction inside the sample. In this case $`\mu _0H/B=1\chi <0`$ which implies thermodynamically unstable sections and the multivaluedness of the induction $`B(H)`$ within some part of each dHvA oscillation period SHOEN ; SOLTBaines . For a long rod-like sample oriented along $`𝐇`$, the instability is avoided by a discontinuous jump $`\delta B=B_2B_1`$ between two stable values $`B_1`$ and $`B_2`$ at a given critical field $`H_c`$. For a plate-like sample perpendicular to $`𝐇`$, the boundary condition $`B=\mu _0H`$ for a uniformly magnetized state leads to domain formation with alternating regions of diamagnetic and paramagnetic magnetization for $`H`$ in the range $`B_1<\mu _0H<B_2`$ CONDON . The proportion of the domains varies with $`H`$ so that $`\overline{B}=\mu _0H`$ is fulfilled as an average over the sample CONWAL ; SOLTBaines . The $`HB`$ diagram is similar to the $`pV`$ diagram of a van der Waals gas, only with more than one discontinuity interval $`\delta B`$, situated periodically on the $`B`$ axis. The existence of domains has been firstly discovered by Condon and Walstedt on a single crystal of silver CONWAL . The domains were revealed by a periodic splitting of the NMR line, corresponding to a local field difference of about 12 G between the paramagnetic and diamagnetic regions in fields of about 9 T. This pioneering result remained the only reference work in the next decades. New experimental possibilities appeared with the development of muon spin rotation. Condon domains were observed in beryllium, white tin, aluminum, lead, and indium SOLT ; SOLTEGOROV . By now, no doubt, Condon domains are expected to appear in pure single crystals of all metals. Thermodynamic aspects of the Condon domain phase transition have been recently treated theoretically GOR98 . While the state of art in this field has been recently reviewed GORDONVagner , some important questions however, concerning domain size and topology, domain wall energy, and pinning properties can only be solved with a detailed knowledge of the domain structure. The state with Condon domains can be considered as physically similar to the intermediate state of type I superconductors, where superconducting and normal regions form in an applied magnetic field. Therefore, domain structures resulting of such different phenomena as superconductivity and dHvA effect may be rather similar. Unfortunately, the magnetic contrast, that is the ratio of $`\delta B`$ to $`B_2`$, is not more than 0.1 $`\%`$ for Condon domains (compared to 100$`\%`$ for the intermediate state). Besides, the magnetic field itself is here hundred times higher. Thus, methods like magnetic decoration or magneto-optical detection used for intermediate state imaging LIVINGSTONE can not be used for Condon domains. In this Letter we present the first experimental results for direct observation of Condon domain structures in silver by a system of ten micro Hall probes being close to the single crystal surface. In the homogeneous state, without domains, all probes show the same dHvA signal $`B(H)`$, i. e. all $`B_i=B(H)`$, where $`i=1,2\mathrm{}10`$ are the Hall probe numbers. In the domain state, the Hall voltages differ between the different probes in the paramagnetic part of the dHvA period. This implies an inhomogeneous magnetic field distribution due to Condon domains at the sample surface. In our measurements the surface of the crystal was either normal to $`𝐇`$ direction or slightly tilted (13). By comparing the data of neighboring Hall probes, new information about the domain structure has been extracted. Fig. 1 shows the Hall probe set-up made of a 1 $`\mu `$m thick Si doped GaAs layer sandwiched between two 10 nm thick undoped GaAs layers. Two arrays of five Hall probes (10$`\times `$10 $`\mu `$m<sup>2</sup> at 40 $`\mu `$m distance) are placed at a distance of $`b=1`$ mm. One array, L, is oriented along the long axis of the sample; the other, T, transverse to this axis. A DC Hall current of 100 $`\mu `$A was applied in series to all five Hall probes of an array. The Hall voltages were read out simultaneously by 5 Keithley multimeters; the arrays L and T were measured one after another. Due to the 3D conducting layer the $`V_i(B)`$ characteristics of the Hall probes were in good approximation linear up to 10 T even at 1.3 K. The correct calibration of the Hall probes was tested at temperatures between 4.2-3.6 K where all Hall probes showed exactly the same dHvA oscillations of the homogenous silver sample. The detection limit of the Hall probes was smaller than 1 G. A high homogeneity (better than 10 ppm in 1 cm<sup>3</sup>) 10 T superconducting magnet was used to set a fixed offset magnetic field $`H_0`$. The slowly varying superimposed field $`H_V`$ ($`\pm 15`$ mT) was made by a watercooled resistive coil. Thus the total applied magnetic field was $`H=H_0+H_V`$. The measurements were performed on a high quality silver single crystal of $`2.4\times 1.6\times 1.0`$ mm<sup>3</sup>. The largest surface of the sample was normal to the $`[100]`$-axis of the crystal. The sample was prepared in the same way as in experiments on radio frequency size effect and time of flight effect (see references in GASPAROV ). The very good quality of the sample results in a very low Dingle temperature, which was estimated from our measurements to be about $`T_D=0.2`$ K. The sample was annealed in O<sub>2</sub> (10<sup>-2</sup> Pa) at 750C during 10 hours. It has a residual resistance ratio $`RRR=R_{300K}/R_{4.2K}=1.6\times `$10<sup>4</sup>, measured by the contactless Zernov-Sharvin method ZERNOVSHARVIN . For a mirror-like surface, the crystal was slightly repolished by 0.1 $`\mu `$m diamond paste after annealing. The surface before polishing had a roughness of about 20-30 $`\mu `$m and no induction splitting due to Condon domains could be observed. The sample was glued by narrow strips of cigarette-paper to the set-up frame to fix the crystal on the Hall probes in order to avoid damage or strain of the single crystal upon cooling down. Moreover, the sample was slightly pressed by a cotton tampon to hold it reliably in high magnetic field. Fig. 2 shows typical $`B(H)`$ traces of Hall probes $`B_1`$ and $`B_5`$ of the L-array over three dHvA periods at 10 T and 1.3 K. In each paramagnetic part of the dHvA period two different inductions are measured at the surface of the sample whereas the induction is homogeneous in the diamagnetic part. The measured traces are reversible for increasing and decreasing magnetic field. We ascribe the measured difference between the induction of neighboring Hall probes to the existence of Condon domains. The maximal induction splitting $`\delta B`$ in a dHvA period was measured as a function of temperature at 10 T (see Fig. 3a) and as a function of field at 1.3 K (see Fig. 3b). At 10 T, the phase boundary is crossed at about 3 K. At 1.3 K, the crossing occurs at about 5 T. The field and temperature range for the occurrence of the induction splitting is in agreement with the $`T`$-$`H`$ phase diagram for the Condon domain state in Ag, as shown in the inset of Fig. 3a for the theoretically calculated phase-diagram of Ag with Dingle temperature 0.2 K GORDON . The solid line in Fig. 3a is the calculated induction splitting $`\delta B`$ GORDON with the phase transition temperature (3 K) and the maximum splitting in silver (12 G) measured by Condon and Walstedt CONWAL as parameters. An anomalous alternating transition order of the Hall probes between the diamagnetic and paramagnetic phase is shown in Fig. 2. Although reproducible, the observed order depends strongly on the experimental configuration. A basically different behavior can be seen between the T- and L-probes in, respectively, Figs. 4 and 5. No regular transition order was observed for T-probes. Sometimes, they transit in ascending (1-5) or in descending order, as if the domain laminae were slightly tilted to the long axis of the sample. Sometimes, as shown in Fig. 4, a middle T-probe transits the last or the first, as if the laminae are bent. In contrast, the order of the L-probes is always 1,2,3,4,5 or reversed as it is shown in Fig. 5a. This implies that the domain structure is approximately laminar with the laminae mainly oriented transverse to the long axis of the sample. However, we found that the L-probe sequence changes alternately between dHvA periods which implies that the domain-wall movement changes direction along the long sample axis between successive dHvA periods. $`\delta B=B_1B_5`$, shown in Fig. 5b, changes sign alternately during four or five periods (see Fig. 6a). Then the, what we will call, ”pendulum” effect breaks down during two periods where the transition order is not clear. After this the pendulum effect repeats. On a slightly tilted sample the pendulum effect disappears completely. Fig. 6 shows the change of the situation after tilting the sample. After rotation around the long sample axis by 13 domains draw up to a regular laminar structure oriented always transverse to the long axis. A similar behavior was observed in white tin in the intermediate state SHARVIN ; LIVINGSTONE indicating the preference of domain walls to align along the sample surface. Furthermore, the transition order is now the same for all dHvA periods (see Fig. 6b). The rotation of the silver single crystal with respect to the magnetic field affects the dHvA frequency spectrum. Only one dHvA frequency (”belly” SHOEN ) remains for the 13 tilted sample. The beating pattern in the oscillatory dHvA signal of the magnetization for the perpendicular field orientation might play a role in the occurrence of the pendulum effect. In this respect we note that the dHvA frequencies (belly orbit at 47300 T and rosette orbit at 19000 T) for the perpendicular field orientation would be compatible with the observed pendulum effect of Fig. 6a. The transitions of the individual Hall probes are very sharp compared to the whole field range of the domain state in neighboring Hall probes (see Fig. 4). This means that the thickness of the domain wall is much smaller than the period of the domain structure. We have never seen more than one transition of a Hall probe in a period. This implies that we saw always only one boundary between para- and diamagnetic phases in an array of 5 successive Hall probes meaning that the period of the domain structure is certainly larger than the distance of $`150\mu `$m between the edge L-probes. This under limit for the domain period ($`p`$) should be compared with the value obtained from the square-root averaged expression $`p\sqrt{wt}`$ for a sample with thickness ($`t`$) and domain wall thickness ($`w`$SHOEN . With the proposed cyclotron radius for the wall thickness ($`1\mu `$m at 10 T in Ag), one obtains at least a 5 times smaller value ($`30\mu `$m) CONWAL . Therefore, from our experiments we find a wall thickness of at least $`20\mu `$m. This is in agreement with the observation that two neighboring middle L-probes at a distance of $`40\mu `$m show often intermediate but different induction values. Therefore, the thickness of a domain wall can not be much smaller than 20 $`\mu `$m. As the real domain pattern turns out to be somewhat bigger than expected, we need either a new set-up with better adapted Hall-probe distances or a scanning Hall probe for more detailed measurements of the domain structure. Exactly the same measurements as presented above were performed on a beryllium sample cut from the same single crystal where Condon domain formation was observed using muon spectroscopy SOLT . The sample was prepared with a surface quality comparable to the Ag crystal. Even though the expected $`\delta B`$ inside the crystal is ten times higher than in silver, we did not find $`\delta B>2`$ G on the sample surface. The attempt of Condon and Walstedt to find domains in beryllium by NMR was not successful, either CONWAL . The authors gave explanations related to the quadrupole broadening and the long nuclear thermalization time in beryllium. However, now we believe that the main reason is the absence of induction splitting $`\delta B`$ at the sample surface. This could be an intrinsic property of beryllium related to its anisotropic magnetostriction LYKOV . In conclusion, Condon domains in silver with induction splitting up to 10 G were observed by micro Hall probes at fields and temperatures which are in agreement with the theoretically estimated phase diagram. A laminar domain structure was found with the orientation mainly transverse to the long sample axis. The domain transitions are always reversible for increasing and decreasing magnet field. For a slightly tilted sample the strange pendulum effect disappears, and the transitions occur in the same order for all dHvA periods. The domain period is not smaller than 150 $`\mu `$m and the domain wall thickness must be about 20 $`\mu `$m. Condon domains in beryllium do not emerge to the surface. We are grateful indebted to M. Schlenker for his support and continuous interest to this work, to I. Sheikin and V. Mineev for fruitful discussions, and to J. Marcus for his help in sample surface preparation. F. Schartner is acknowledged for the preparation of the Hall probes.
warning/0506/gr-qc0506068.html	ar5iv	text	# Is empty spacetime a physical thing? ## 1 Introduction Since its inception, the ether has proved a troubled notion which, contrary to common belief, still haunts physics at the turn of a new century. What can one learn from the history of this ether and in what sense is it still of utmost concern to physics? The need for an ether as a material medium with mechanical properties first became apparent to Descartes in the first half of the seventeenth century in an attempt to avoid any actions that would propagate, through nothing, from a distance. During its history of roughly three centuries in its original conception, the idea of the ether managed to materialize in endless forms via the works of countless investigators. Its main purpose of providing a medium through which interactions could propagate remained untouched, but the actual properties with which it was endowed in order to account for and unify an ever-increasing range of phenomena were mutually incongruent and dissimilar. Never yielding to observation and constantly confronted by gruelling difficulties, the mechanical ether had to reinvent itself continually, but its very notion staggered not a bit. After almost 300 years of bitter struggle, the mechanical ether eventually gave in. The first step of this change took place in the hands of Lorentz, for whom the ether was a sort of substantial medium that affected bodies moving through it not mechanically but only dynamically, i.e. due to the fact that bodies moved through it. Drude and Larmor further declared that the ether need not actually be substantial at all but simply space with physical properties. The second and decisive step in this de-mechanization of the ether was brought about by Poincaré and Einstein through ideas that eventually took the shape of the special relativity theory. In Einstein’s (1983) own words, this change “consisted in taking away from the ether its last mechanical quality, namely, its immobility” (p. 11). Far from being dead, however, the ether had only transmuted its character—from a mechanical substance to an absolute inertial spacetime. The need for regarding inertial spacetime as an ether came after noticing that “empty” spacetime, despite being unobservable and unalterable, displayed physical properties, such as providing a reference for acceleration via its geodesics. The nature of this new ether underwent yet another change with the theory of general relativity. According to Einstein (1961, p. 176), the ether was now spacetime’s dynamic and intrinsic metric content. This was so significant a change that it modified the very ideas of ether and empty spacetime. By making metric spacetime alterable, it actually put an end to its status as a genuine ether. And by making the metric field a content of spacetime, it did away with empty spacetime, since now to vacate spacetime means to be left with nothing at all. In fact, one is left with something: the spacetime points; however, these had been denied physical reality by Einstein’s hole argument. From this standpoint, Einstein concluded that empty spacetime cannot possess any physical properties, i.e. that empty spacetime does not exist. It is the purpose of this article to challenge the certainty of this conclusion. The way to achieve this will be connected with the main ideas presented in our previous article (Meschini, Lehto, & Piilonen, 2005), where the need to study the problem of the nature of empty spacetime (i) equipped with quantum theory and some guiding physical principle directly relevant to its existence—here proposed to be that of diffeomorphism invariance—and (ii) from a non-geometric point of view was put forward. The first requirement stemmed from the simple fact that, without an appropriate guiding light, the search for new ideas becomes pure guesswork; the second originated from the observation that any genuinely new understanding of empty spacetime may require going beyond its geometric characterization entirely (further reasons justifying this second requirement will be given in this article). In particular both these observations were turned into criticisms of what is currently understood as pregeometry—there deemed a considerable incongruity. When further heeding the history of the ether, the question of the physical existence of empty spacetime must also involve that observables intrinsic to empty spacetime be found so as to be able to support any sound claim as to its reality. The present article constitutes the beginning of this endeavour. ## 2 The mechanical ether It is appropriate to start this investigation by tracing the rich history of the ether, starting here from its older conception as a material substance, passing through its virtual disappearance after the progressive removal of its mechanical attributes, and ending with its new form of an immaterial, geometric substratum, as analyzed in the next two sections. For the historical review of this section, the very comprehensive work of Whittaker (1951) will be followed as a guideline.<sup>2</sup><sup>2</sup>2Page numbers in parentheses in this section are to Vol. 1 of this reference unless otherwise stated. René Descartes (1596–1650) was the first to introduce the conception of an ether as a mechanical medium. Given his belief that action could only be transmitted by means of pressure and impact, he considered that the effects at a distance between bodies could only be explained by assuming the existence of a medium filling up space—an ether. He gave thus a new meaning to this name, which in its original Greek ($`\alpha \stackrel{`}{\iota }\theta \stackrel{´}{\eta }\rho `$) had meant the blue sky or the upper air. The ether was unobservable and yet it was needed to account for Descartes’ mechanistic view of the universe, given his said assumptions.<sup>3</sup><sup>3</sup>3With his invention of the coordinate system, Descartes was, at the same time, the unwitting precursor of the later conception of the ether as a form of space. At the same time, the notion of an ether was right from its inception entwined with considerations about the theory of light. Descartes himself explained the propagation of light as a transmission of pressure from a first type of matter to be found in vortices around stars to a second type of matter, that of which he believed the ether to be constituted. (pp. 5–9) The history of the ether continued tied to the theory of light with Robert Hooke’s (1635–1703) work. In an improvement with respect to Descartes’ view, he conceived of light as a wave motion, an exceedingly small vibration of luminous bodies that propagated through a transparent and homogeneous ether in a spherical manner. Hooke also introduced thus the fruitful idea of a wave-front. (pp. 14–15) Isaac Newton (1642–1727) rejected Hooke’s wave theory of light on the grounds that it could not explain the rectilinear propagation of light or its polarization (see below). In its place, Newton proposed that light consisted of rays that interacted with the ether to produce the phenomena of reflection, refraction and diffraction, but that did not depend on it for their propagation. He gave several options as to what the true nature of light might be, one of which was that it consisted of particles—a view that later on would be associated with his name; nevertheless, as to the nature of light, he “let every man here take his fancy.” Newton also considered it possible for the ether to consist of different “ethereal spirits,” each separately suited for the propagation of a different interaction. (pp. 18–20) Regarding gravitation in the context of his universal law of attraction, Newton did not want to pronounce himself as to its nature. He nonetheless conjectured that it would be absurd to suppose that gravitational effects could propagate without the mediation of an ether. However, Newton’s eighteenth century followers gave a twist to his views; antagonizing with Cartesians due to their rejection of Newton’s gravitational law, they went as far as denying the existence of the ether—originally Descartes’ concept—and attempted to account for all contact interactions as actions at a distance between particles. (pp. 30–31) Christiaan Huygens (1629–1695) was also a supporter of the wave theory of light after observing that light rays that cross each other do not interact with one another as would be expected of them if they were particles. Like Hooke, he also believed that light consisted of waves propagating in an ether that penetrated vacuum and matter alike. He managed to explain reflection and refraction by means of the principle that carries his name, which introduced the concept of a wave front as the emitter of secondary waves. As to gravitation, Huygens’ idea of a Cartesian ether led him to account for it as a vortex around the Earth. (pp. 23–28) An actual observation that would later have a bearing on the notions of the nature of light and of the ether was that of Huygens’ regarding the polarization of light. He observed that light refracted once through a so-called Icelandic crystal, when refracted through a second such crystal, could or could not be seen depending on the orientation of the latter. Newton correctly understood this result as the first light ray being polarized, i.e. having properties dependent on the directions perpendicular to its direction of propagation. He then concluded that this was incompatible with light being a (longitudinal) wave, which could not carry such properties. (pp. 27–28) Another thoroughly Cartesian account of the ether was presented by John Bernoulli (1710–1790), Jr., in an attempt to provide a mechanical basis for his father’s ideas on the refraction of light. Bernoulli’s ether consisted of tiny whirlpools and was interspersed with solid corpuscules that could never astray much from their average locations. A source of light would temporarily condense the whirlpools nearest to it, diminishing thus their centrifugal effects and displacing the said corpuscules; in this manner, a longitudinal wave would be started. (pp. 95–96) In the midst of a general acceptance of the corpuscular theory of light in the eighteenth century, also Leonhard Euler (1707–1783) supported the view of an ether in connection with a wave theory of light after noticing that light could not consist of the emission of particles from a source since no diminution of mass was observed. Most remarkably, Euler suggested that, in fact, the same ether served as a medium for both electricity and light, hinting for the first time at a unification of these two phenomena. Finally, he also attempted to explain gravitation in terms of the ether, which he assumed to have more pressure the farther from the Earth, so that the resulting net balance of ether pressure on a body would push it towards the centre of the Earth. (pp. 97–99) At the turn of the century, the wave theory of light received new support in the hands of Thomas Young (1773–1829). Within this theory, Young explained reflection and refraction in a more natural manner than the corpuscular theory and, more importantly, he also accounted successfully for the phenomena of Newton’s rings (and hinted at the cause of diffraction) by introducing an interference principle for light waves. It was Augustin Fresnel (1788–1827) who, in 1816 and amidst an atmosphere of hostility towards the wave theory, managed to explain diffraction in terms of Huygens’ and Young’s earlier findings. (pp. 100–108) Young and Fresnel also provided an alternative explanation of stellar aberration, which had been first observed by James Bradley (1692–1762) in 1728 while searching to measure stellar parallax, and which had so far been explained in terms of the corpuscular theory of light. Young first proposed that such effect could be explained assuming that the Earth did not drag the ether with it, so that the Earth’s motion with respect to it was the cause of aberration. Subsequently, Fresnel provided a fuller explanation that could also account for aberration being the same when observed through refractive media. Following Young, Fresnel suggested that material media partially dragged along the ether in such a way that the latter would pick a fraction $`11/n^2`$ (where $`n`$ is the medium’s refractive index) of the medium’s velocity. So far the ether was viewed as a somewhat non-viscous fluid that could be dragged along in the inside of refractive media in proportion to their refractive index, and whose longitudinal excitations described light. (pp. 108–113) Considerations about the polarization of light would bring along fundamental changes to the conception of the ether. As Newton had previously observed, the properties of polarized light did not favour a longitudinal-wave theory of light. Inspired by the results of an experiment performed by François Arago (1786–1853) and Fresnel, Young hit on the solution to the problem of polarization by proposing that light was a transverse wave propagating in a medium. Fresnel further hypothesized that the ether must then be akin to a solid and display rigidity so as to sustain such waves. (pp. 114–117) The fact that only a rigid ether could support transverse waves robbed the idea of an immobile, undragged ether of much of its plausibility, since it was hard to imagine a solid medium of some sort that would not be, at least, partially dragged by bodies moving through it. George Stokes (1819–1903) rose up to this challenge by providing a picture of the ether as a medium that behaved like a solid for high-frequency waves and as a fluid for slow moving bodies. As a fluid, Stokes’ ether was dragged along by material bodies such that, in particular, it was at rest relative to the Earth’s surface. (pp. 128, 386–387) Michael Faraday (1791–1867) gave a new dimension to the ether conception by introducing the notion of field, which in hindsight was the most important concept to be invented in this connection.<sup>4</sup><sup>4</sup>4And this not without a sense of irony. This is so because, at first, its was on the field, a physical entity existing on its own and needing no medium to propagate, that the overthrow of the mechanical ether would rest; however, later on Einstein would reinstate the ether as a (metric) field itself. In his studies of the induction of currents, of the relation of electricity and chemistry, and of polarization in insulators, he put forward the concepts of magnetic and electric lines of force permeating space. He introduced thus the concept of a field as a stress in the ether and present where its effects took place. He went on to suggest that an ether may not be needed if one were to think of these lines of force—themselves part of material bodies—as the carriers of transverse vibrations, including light and radiant heat as well. Or then that, if there was an luminiferous ether, it might also carry magnetic force and “should have other uses than simply the conveyance of radiations.” By including also the magnetic field as being carried by the ether, Faraday added to Euler’s earlier prophecy, and he hinted for the first time at the conception of light as an electromagnetic wave. (pp. 170–197) Another unifying association of this type was made by William Thomson (1824–1907), who in 1846 compared heat and electricity in that the isotherms of the former corresponded to the equipotentials of the latter. He suggested furthermore that electric and magnetic forces might propagate as elastic displacements in a solid. James Clerk Maxwell (1831–1879), inspired by Faraday’s and W. Thomson’s ideas, strove to make a mechanical picture of the electromagnetic field by identifying static fields with displacements of the ether (for him equivalent to displacements of material media) and currents with their variations. At the same time, Maxwell, like Gustav Kirchhoff (1824–1887) before him, was impressed by the equality of the measured velocity of light and that of electromagnetic disturbances of his theory, and suggested that light and electromagnetic waves must be waves of the same medium. (pp. 242–254) So far, the theories of Maxwell and Heinrich Hertz (1857–1894) had not made any distinction between ether and matter, with the former considered as totally carried along by the latter. These theories were still in disagreement with Fresnel’s successful explanation of aberration in moving refractive media, which postulated a partial ether drag by such bodies. However, experiments to detect any motion of the Earth with respect to the ether, such as those by Albert Michelson (1852–1931) and Edward Morley (1838–1923), had been negative and lent support to Stokes’ theory of an ether totally dragged at the surface of the Earth. (pp. 386–392) Not content with Stokes’ picture, in 1892 Hendrik Lorentz (1853–1928) proposed an alternative explanation with his theory of electrons, which reconciled electromagnetic theory with Fresnel’s law. However, Lorentz’s picture of the ether was that of an electron-populated medium whose parts were mutually at rest; Fresnel’s partial drag was therefore not allowed by it. Lorentz’s theory denied the ether mechanical properties and considered it only space with dynamic properties (i.e. affecting bodies because they moved through it), although still endowed it with a degree of substantiality.<sup>5</sup><sup>5</sup>5See (Kostro, 2000, p. 18) and (Kox, 1989, pp. 201, 207). The negative results of the Michelson-Morley experiment were then explained by Lorentz by means of the existence of FitzGerald’s contraction, which consisted in a shortening of material bodies by a fraction $`v^2/2c^2`$ of their lengths in the direction of motion relative to the ether. Thus, the ether would cease being a mechanical medium to become a sort of substantial, dynamic space. (pp. 392–405) Near the end of the nineteenth century, the conception of the ether would complete the turn initiated by Lorentz with the views of Paul Drude (1863–1906) and Joseph Larmor (1857–1942), which entirely took away from the ether its substantiality. Drude (1894, p. 9) declared: > Just as one can attribute to a specific medium, which fills space everywhere, the role of intermediary of the action of forces, one could do without it and attribute to space itself those physical characteristics which are now attributed to the ether.<sup>6</sup><sup>6</sup>6Quoted from (Kostro, 2000, p. 20). Also Larmor claimed that the ether should be conceived as an immaterial medium, not a mechanical one; that one should not attempt to explain the dynamic relations so far found in terms of > mechanical consequences of concealed structure in that medium; we should rather rest satisfied with having attained to their exact dynamical correlation, just as geometry explores or correlates, without explaining, the descriptive and metric properties of space.<sup>7</sup><sup>7</sup>7Quoted from (Whittaker, 1951, Vol. 1, p. 303). Larmor’s statement is so remarkable that it will receive more attention later on in Section 5. Despite the seeming superfluousness of the ether even taken as a fixed dynamic space, Lorentz held fast to the ether until his death, hoping perhaps that motion relative to it could still somehow be detected (Kox, 1989). Others, like Poincaré and Einstein, understood the repeated failed attempts to measure velocities with respect to the ether as a clear sign that the ether, in fact, did not exist. Henri Poincaré (1854–1912) was the first to reach such a conclusion; in 1899 he asserted that absolute motion with respect to the ether was undetectable by any means, and that optical experiments depended only on the relative motions of bodies; in 1900 he openly distrusted the existence of the ether with the words “Our ether, does it really exist?”; and in 1904 he proposed a principle of relativity. (Vol. 2, pp. 30–31) It was Albert Einstein (1879–1955) who in 1905 provided a theory where he reinstated these earlier claims but with a new, lucid interpretational basis. In particular, Einstein considered > \[T\]he introduction of a “luminiferous ether”…to be superfluous inasmuch as the view here to be developed will not require an “absolutely stationary space” provided with special properties…(Einstein, 1952a, p. 38). Thus, in the hands of Poincaré and Einstein, the ether had died. ## 3 Einstein’s ethers Belief in the non-existence of the ether would, nevertheless, not last very long, for it would soon rise from its ashes—transmuted. A rebirth of the ether was now advocated by Einstein on the grounds that, without it, “empty” space could not have any physical properties; yet it displayed them through the effects of absolute acceleration. It is a well-known fact that all motion simply cannot be reduced to the symmetric relationship between any two reference systems, as Newton’s (1962, pp. 10–12) rotating-bucket and revolving-globes (thought) experiments, and Einstein’s (1952b, pp. 112–113) rotating-spheres thought experiment aimed to show. Some reference frames clearly show their privileges. In Newtonian mechanics and in special relativity, these are the inertial frames; in general relativity, these are the freely-falling frames. Such directly unobservable frames of reference confer physical properties on “empty” spacetime, and were held by Einstein as a new embodiment of the ether. In order to understand Einstein’s conceptions,<sup>8</sup><sup>8</sup>8See (Kostro, 2000) for a useful source of material on Einstein and the ether. However, note that this book does not deal with the issue of the hole argument and the reality of spacetime points. we now develop a further, tentative characterization of the ether by distinguishing three different, possible realizations of it. * An ether is an entity that has sources and that cannot be observed directly, although it can be observed indirectly. This is the case of physical fields such as the metric, electric and magnetic fields, $`𝐠(x)`$, $`\stackrel{}{E}(x)`$ and $`\stackrel{}{B}(x)`$, which can be detected through the behaviour of test particles and influenced through changes in their sources. We name these ethers electric, magnetic, gravitational, etc. ether, as the case might be. * An ether is an entity that does not have sources—therefore cannot be influenced in any way—and cannot be observed directly, although it can be observed indirectly via the behaviour of test particles. Two constructs that realize this notion are Newtonian space and special relativistic spacetime (characterized by the constant metric field $`𝜼`$). We name these ethers inertial ethers. * An ether is an entity that does not have sources—therefore cannot be acted upon in any way—and that “acts” but cannot be observed directly, nor indirectly through its effects on test particles. This kind of ether is realized by the spacetime points, about which much will be said below and in Sections 4 and 5. For reasons that will become clear shortly, we name this ether a geometric ether. Einstein called Newtonian space and the metrics of special and general relativity ethers. What reasons lie behind these identifications? The Newtonian and special relativistic inertial ethers above are ethers as genuine as the mechanical one had once been. By this we mean that, after allowing for a conceptual change from mechanical medium to geometric space, these inertial ethers are also directly unobservable substrata that nonetheless are thought to cause distinct effects in the observable world. While the mechanical ether—in its endless varieties—was regarded as the cause of effects ranging from gravitational to caloric, the inertial ethers afforded an explanation of motion. In effect, born out by his experiments according to which (the effect of) acceleration was not relative and rotation in empty space was meaningful, Newton held on to the need for a space absolutely at rest with respect to which this absolute acceleration could be properly defined.<sup>9</sup><sup>9</sup>9In fact, only a family of inertial spaces linked by Galilean transformations would have sufficed judging from the properties of the theory itself. Newton’s absolute space is therefore an “old new ether” in sense (ii), since it has effects on test particles but has no sources and cannot be affected in any way. From the point of view of the theory of special relativity, the inertial ether is needed for the same reasons as above, but it cannot be assigned any mechanical property whatsoever—not even immobility or rest; any talk about its state of motion is meaningless (Einstein, 1983, p. 13). This ether of special relativity is nothing other than a background inertial spacetime, i.e. an infinite family of inertial frames linked by Lorentzian transformations. It is characterized by the everywhere-constant metric field $`𝜼`$ and represents an indirectly observable “empty” spacetime endowed with physical properties: it defines the standards of space, time and motion for a test particle in an otherwise empty world. In spite of the historical success of the inertial ethers, we need not concern ourselves with them in this investigation. The reason is simply that, as independent concepts, they were overthrown by the theory of general relativity. Indeed, according to our best available understanding, the inertial ether is not an accurate description of Nature but only applies in certain limiting situations, and can be understood as a special case of the general relativistic metric $`𝐠(x)`$. General relativity conferred a totally new meaning on the concept of ether—so much so that the very notion of ether cannot be applied to metric spacetime any longer. Like before, the physical properties of spacetime are carried by its metric field—now $`𝐠(x)`$—which again defines the standards of space, time and motion. However, this metric field is a dynamic magnitude that is subject to change as dictated by the distribution of matter. In other words, the ether is no longer immutable but is revealed to have matter as its source, through which it can be acted upon. Einstein had been explicitly concerned with this problem of finding an influenceable field to replace the prevailing immovable $`𝜼`$. When dealing with his rotating-spheres thought experiment, he said: > What is the reason for this difference \[spherical and ellipsoidal\] in the two bodies? No answer can be admitted as epistemologically satisfactory, unless the reason given is an observable fact of experience…\[T\]he privileged space $`R_1`$ \[inertial\] of Galileo…is merely a factitious cause, and not a thing that can be observed…The cause must therefore lie outside this system…\[T\]he distant masses and their motions relative to $`S_1`$ and $`S_2`$ \[the spheres\] must then be regarded as the seat of the causes (which must be susceptible to observation) of the different behaviour of our two bodies $`S_1`$ and $`S_2`$. They take over the role of the factitious cause $`R_1`$. (Einstein, 1952b, pp. 112–113) This shows that Einstein was looking for sources via which the physical properties of spacetime could be influenced, and thus no longer fixed and beyond reach. In this sense, Einstein was the first physicist since the conception of the ether to bring it to full physical accountability (i.e. not only to passively observe its effects but also to control its structure) by finding, through the field equation, the ether’s intimate linkage to matter as its source. He achieved this, however, not without altering the meaning of the original concept. To be sure, this new property of $`𝐠(x)`$ greatly upsets its interpretation as an ether, since it puts it on virtually the same footing as other physical fields. Indeed, the metric field $`𝐠(x)`$ is akin to the electric field $`\stackrel{}{E}(x)`$ or magnetic field $`\stackrel{}{B}(x)`$ (or to the electromagnetic tensor field $`𝐅(x)`$). Whereas matter acts as the source of the former, charges act as sources of the latter; further, e.g. like seeing a compass move is evidence of the existence of $`\stackrel{}{B}(x)`$, so is seeing a stone fall evidence of $`𝐠(x)`$. The presence of all these fields can thus be granted beyond reasonable doubt.<sup>10</sup><sup>10</sup>10And yet, this comparison is not as straightforward as one might wish. While electromagnetic fields carry physical units, the metric field does not, and neither does it result as the dimensionless quotient of other magnitudes. This makes the indirect observability of the metric field more intricate in comparison. It might be argued that one observes the effects of spacetime curvature and not really of the metric field. Then, in order to understand Einstein’s views regarding $`𝐠(x)`$ as an ether, we must recognize that the word can be meant not only in its later, negative sense of something physically unreal but also in its more primitive sense of an underlying, ubiquitous substratum with physical properties—even if now it were on a par with other physical fields. While Einstein’s usage highlights the all-pervading character of an ether, it blurs the more important issue of its physical reality or non-reality. As we saw above, as far as observability is concerned, the metric is a physical field analogous to the electromagnetic field; however, while the former necessarily permeates all of spacetime (a ubiquitous substratum), the latter does not. Like Einstein, also Weyl (1918, p. 182) recognized this point early on: > The coefficients of the fundamental metric form are therefore not simply the potentials of the gravitational and centrifugal forces, but determine in general which points of the universe are in reciprocal interaction. For this reason the name “gravitational field” is perhaps too unilateral for the reality described by this expression and should better be replaced by the word “ether;” while the electromagnetic field should simply be called field.<sup>11</sup><sup>11</sup>11Quoted from (Kostro, 2000, p. 74). In order to also emphasize the aspect of the issue relating to physical existence, we reserve the use of the word “ether” in both of the said senses, i.e. for denoting an underlying substratum that acts (has physical properties) but cannot be acted upon or observed through physical tests. For one or more reasons given above, this description is not realized by the electromagnetic field, but neither by the general-relativistic metric. However, the description is realized by the spacetime points, since these geometric objects “perform the localization of fields” but they cannot be observed or influenced themselves as expressed by the diffeomorphism invariance postulate of general relativity (see Section 4). Therefore, we call spacetime points the geometric ether.<sup>12</sup><sup>12</sup>12Notice that this denomination would be too inclusive if we had also accepted the metric field $`𝐠(x)`$ as an ether (or had concerned ourselves with $`𝜼`$ as one), since its character is certainly also geometric. In any case, $`𝐠(x)`$ would have been a quantitatively geometric ether (geometric magnitude), whereas spacetime points are only a qualitatively geometric ether (geometric object). Finally and in the same connection, not only did general relativity change the very notion of ether, but also that of empty spacetime. It is another feature of $`𝐠(x)`$ that it depends on the spacetime coordinates (geometrically speaking, on its points), so that it cannot consist of an absolute background associated with empty spacetime as $`𝜼`$ in special relativity. On the contrary, the metric field $`𝐠(x)`$ constitutes an intrinsic content<sup>13</sup><sup>13</sup>13Einstein (1961) distinguished space from its contents with the words: “\[S\]pace as opposed to ‘what fills space,’ which is dependent on the coordinates…” (p. 176) \[Italics added\]. This distinction makes sense since anything that depends on the spacetime coordinates (or points) must be in spacetime. of spacetime. Its removal means that the standards of space, time and motion (i.e. spacetime’s geometric structure) are gone, so that not empty spacetime but, rather, nothing remains without it (Einstein, 1961, p. 176). Einstein’s conclusion will be amplified in Section 4, where the physical status of what does remain—the spacetime points—will be investigated further. ## 4 The geometric ether Einstein’s conclusion that spacetime, empty of its metric-field “ether,” does not exist demands a qualification. This can be seen as soon as one realizes that, after ridding spacetime of its geometry, it is not “nothing” which remains but the spacetime points themselves. However, when Einstein (1961, 1983) seemingly jumped to conclusions in these (and other) expositions, he might have already been assuming as known the results of his so-called hole argument: spacetime points are not physical either, so that they could not constitute truly empty spacetime in any physical sense. Given that Einstein’s field equation is generally covariant, if $`𝐠(x)`$ is a solution to it corresponding to the matter-content source $`𝐓(x)`$, then so is $`𝐠^{}(x^{})`$ with corresponding matter-content source $`𝐓^{}(x^{})`$, for any continuous coordinate transformation $`xx^{}`$. This is simply so because the same matter and metric fields are being viewed from two different frames of reference, both being equivalent for the description of Nature. Dropping the primed frame completely, how to express $`𝐓^{}(x^{})`$ and $`𝐠^{}(x^{})`$ as viewed from the unprimed frame? All that needs to be done is replace $`x^{}`$ by $`x`$ so that the matter and metric content are now seen from $`S`$ as earlier seen from $`S^{}`$. This transformation consists of the active equivalent of a transformation of coordinates (from one reference frame to any another), and achieves its goal by “moving” the spacetime points and not the frame of reference. The result is that, since all frames of reference are physically equivalent, both $`𝐠(x)`$ with source $`𝐓(x)`$, and $`𝐠^{}(x)`$ with source $`𝐓^{}(x)`$ are solutions to the field equation as seen from the same frame. Now let $`H`$ be a so-called hole<sup>14</sup><sup>14</sup>14The hole argument has been reviewed in innumerable places. See e.g. (Stachel, 1989, pp. 71–81), where relevant quotations and a full list of early references can be found. in spacetime in the sense that within this region no matter content is present, i.e. $`𝐓(x)=0`$; outside $`H`$, on the other hand, $`𝐓(x)`$ is non-null. Furthermore, let $`𝐠(x)`$ be the metric content of spacetime, necessarily non-null both inside and outside the hole, and $`\varphi `$ an active coordinate transformation with the property that it is equal to the identity outside $`H`$ and different from the identity inside $`H`$; demand also that the transformation be continuous at the boundary of $`H`$. The property of general covariance of Einstein’s field equation now implies that the unchanged $`𝐠(x)`$, with unchanged source $`𝐓(x)`$, is a solution outside $`H`$; inside $`H`$, where no matter content is present, both $`𝐠(x)`$ and $`𝐠^{}(x)`$ are (mathematically distinct) solutions. Einstein’s conclusion was that, their mathematical differences notwithstanding, $`𝐠(x)`$ and $`𝐠^{}(x)`$ must be physically the same or else the field equation would not be causal, since both metric fields are produced by the same source outside $`H`$. The only way to avoid this dire consequence was for Einstein to postulate the unreality of spacetime points since, in this manner, the above active transformation of the metric field does not entail that there should be any observational differences.<sup>15</sup><sup>15</sup>15Notice that only the hole argument applied to general relativity leads to the physical unreality of spacetime points. In special relativity, for example, the absolutely given metric field $`𝜼`$ can be used to set up an inertial frame with respect to which spacetime points become physical events. In general relativity, on the other hand, there is no such possibility until the metric field $`𝐠(x)`$ has been obtained as a solution to Einstein’s field equation; however, herein lies the problem: this field equation is not causal unless spacetime points are held to be physically unreal. See (Stachel, 1989, p. 78). There would be nothing remarkable about this conclusion if spacetime points were otiose objects that contributed nothing to our physical understanding. Yet that is not the case. The concept of field as a localized physical magnitude (whether scalar, vectorial or tensorial), on which so much of our scientific portrayal of the world (not just spacetime theories) is based, rests on the notion of point in order to have any meaning. In other words, points are intrinsic to the very concept of physical field. In this respect, Auyang wrote: > The spatiotemporal structure is an integral aspect of the field…We can theoretically abstract it and think about it while ignoring the dynamical aspect of the field, but our thinking does not create things of independent existence…(Auyang, 2001, p. 214) He further criticized Earman’s remarks that > When relativity theory banished the ether, the space-time manifold $`M`$ began to function as a kind of dematerialized ether needed to support fields. (Earman, 1989, p. 155) We agree with Auyang’s view that “spacetime” (spacetime points) is not merely a substratum on which to mathematically define fields. Spacetime points are inherent in fields inasmuch as they perform the physical task of localizing the latter. This is what, after Weyl (1949), Auyang (2001, p. 209) called the “this” or “here-now” aspect of a field, additional to its “thus” or qualitative aspect. In this view, spacetime points satisfy the traditional label of “unobserved actors not acted upon,” and hence our interpretation of points as an ether. However, although we also concur with Auyang in that points are the illusive creations of our geometric thinking, we do not offhand renounce the possibility that empty spacetime—beyond its geometric description—be real on its own, and not simply an illusion of our brains: might it not still be observable as a stand-alone entity? A case in favour of the physical reality of spacetime points, despite the odds against them, was put forward by Friedman (1983, pp. 216–263). According to him, the core of an explanation of natural phenomena is to be able to reduce a wide variety of them to a single framework (Friedman, 1974), so that what one is required to believe is not a vast range of isolated representational structures, but a single, all-encompassing construct. However, in order to provide scientific understanding, theoretical entities with sufficient unifying power must be taken to be of a literal kind. Thus, Friedman argued that denying the existence of spacetime points—themselves essential for geodesics to exist—can only lead to a loss of unifying and explanatory power in spacetime theories. As appealing as this argument may be, it is defeated by a look at history itself. If anything has been learnt from the narration of the history of the old ether, it should be this: no other conception gripped the minds of so many illustrious men of science for longer and more strongly than that of the ether. It was always held to be physically real and unified a wide range of phenomena (light, heat, gravitation, electricity and magnetism) despite its relentlessly unobservable existence. And yet, at the opening of the twentieth century, it became superfluous and useless, and was declared nonexistent. Evidently, no spell of time during which a conception proves to be extremely successful is long enough to declare it real because of its utility or unifying power. What is to guarantee that today’s extremely successful spacetime points—our geometric ether—will not run, in their own due time, the same fate as their mechanical ancestor? If any concept that is of aid in physics is to be held real, nothing more and nothing less should be demanded of it that it be observable. Now, in studying current theoretical pictures comprised under the generic label of quantum gravity (including pregeometry), one can notice that their overall trend is not to try to overcome the geometric ether as here explained by attempts to observe empty spacetime, but to create ever more involved metric (i.e. non-empty) instantiations of it and then, possibly, to observe those. In other words, quantum gravity has come to be an attempt to replace general relativity’s quantitatively geometric picture of spacetime by other equally quantitatively geometric pictures which, at the same time, include unobservable geometric objects. In this respect, Brans (1999, pp. 597–602) argued that much by way of directly unobservable structure is taken for granted in current spacetime theories, such as the existence of a point set, a topology, smoothness, a metric, etc. He compared the vortices and atoms of the mechanical ether with these unobservable building blocks of spacetime and, moreover, with the yet considerably more complex spacetime structures and “superstructures” devised more recently. We interpret Brans to be asking: will strings and membranes, spin networks and foams, nodes and links—to name but a few—appear a hundred years from now like the mechanical ether does today? Are the above the new geometric counterparts of the old mechanical contraptions? Be this as it may, even if these more recently proposed structures were not to follow such a fate and could eventually be observed, the observations relating to them would give evidence of a new quantitatively geometric constitution of spacetime—therefore, of a filled spacetime.<sup>16</sup><sup>16</sup>16Hypothetical observation of the effects of a spacetime lattice (e.g. Smolin, 2004, p. 64) could perhaps be argued as indirect evidence of the vertices themselves, but certainly also as the effects of the lattice’s full metric structure. However, if one’s endeavour is to advance the understanding of empty spacetime, then one must move away from matter and geometry—forwards on a new path. ## 5 Beyond the geometric ether Is Einstein’s conclusion (Section 4) that there is nothing beyond the gravitational “ether” $`𝐠(x)`$ final, then? It is not as long as an issue remains. Might it not be possible to observe this “nothing,” the geometric ether? Can one find observables intrinsic to empty spacetime itself? The story of the mechanical ether is repeating itself today in a geometric, instead of mechanical, guise. The old ether was superseded by stripping it of all its intrinsic, mechanical properties and rendering them superfluous; the old ether did not exist. The new ether is also unobservable as the old one and as problematic, in the sense that it hints at the presence of an unresolved physical issue: from where observable does empty spacetime’s physical capacity to “localize fields” come? To answer this question, the new ether must, like its ancestor, have another layer of its nature revealed, again by stripping it of its intrinsic properties: this time, geometric ones. Our belief is that, having gone beyond the geometric ether, i.e. beyond geometry entirely, observables intrinsic to empty spacetime might be identified. To put it differently, we are here proposing an updated version of Larmor’s centenary words: > We should not be tempted towards explaining the simple group of relations which have been found to define the activity of the aether by treating them as mechanical consequences of concealed structure in that medium; we should rather rest satisfied with having attained to their exact dynamical correlation, just as geometry explores or correlates, without explaining, the descriptive and metric properties of space.<sup>17</sup><sup>17</sup>17Quoted from (Whittaker, 1951, Vol. 1, p. 303). Larmor’s statement is remarkable for its correctness, and all the more remarkable for its incorrectness. Its first half (up to the semicolon) is a correct testimony of what soon would prove itself the way out of the mechanical ether problem: its denial by special relativity. Its second half was—at the time of its utterance—a correct comparison between the way Larmor thought the mechanical ether should be considered and the way geometry was then regarded, namely, not as background for the explanation of phenomena. However, 16 years later in 1916 the ether would be overtly<sup>18</sup><sup>18</sup>18In fact, Newton’s space had already assumed this role over 200 hundred years before, and special relativity’s inertial spacetime 11 years before (see Section 3), but neither of them had been openly considered as spaces with physical geometric properties until after 1916. reinstated as a dynamic field, and the physical properties of spacetime would be explained by it in the sense of a geometric substratum. This trend of attempting to explain phenomena in terms of geometry did not stop with general relativity but continued with the efforts of quantum gravity and, in particular, of so-called pregeometry. Therefore, nowadays geometry does explain, as a substratum, the properties of spacetime, and the second part of Larmor’s view is no longer correct. The updated version of Larmor’s words proposed here reads thus—We should not be tempted towards explaining the simple group of relations (fields’ local aspect) which have been found to define the activity of the geometric ether (spacetime points) by treating them as consequences of concealed structure in that geometric medium; we should rather seek to explain the activity of the geometric ether beyond its geometric nature, searching for observables intrinsic to empty spacetime via non-geometric concepts. Unlike Larmor, we do not renounce an explanation of the geometric ether, but we do not attempt to find it at the same conceptual level as this ether finds itself. Our search for spacetime observables requires taking a conceptual step beyond the state of affairs as left by Einstein’s hole argument. We believe that the current philosophical literature on this problem (e.g. Butterfield, 1989; Earman & Norton, 1987; Rynasiewicz, 1994, 1996) has not been able to take this step by adding something physically new to the discussion. On the contrary, the philosophical debate appears to function in the spirit of Earman’s (1989) words, which measure the fruitfulness of a work by asking “How many discussions does it engender?” (p. xi). What is needed is a new physical insight by means of which the present philosophical debate may be rendered inconsequential, much like the older disputes as to the shape and position of the Earth or the nature of the heavenly bodies were only settled by new physical investigations. In order to understand how our step beyond the hole argument is related to the hole argument itself, we will now re-rehearse it but this time in geometric, rather than coordinate, language and for a general field. Let there be two points, $`P`$ and $`Q`$, inside $`H`$ linked by a diffeomorphism $`\varphi (P)=Q`$, and let each of them be the local aspect, or location, of fields $`f(P)`$ and $`f^{}(Q)`$, respectively, solutions to the field equation with the same source. The demand that $$f(P)=f^{}(Q)$$ (1) (cf. $`𝐠(x)=𝐠^{}(x^{})`$) reproduces geometrically the requirement that, after a coordinate transformation $`xx^{}`$, a physical situation remains unchanged. In terms of an active transformation, point $`P`$ (cf. $`x`$) is dropped, and $`f(P)`$ is now represented only in terms of $`Q`$ (cf. $`x^{}`$) by its push-forward $$\varphi ^{}[f(P)]=f[\varphi (P)]=f(Q).$$ (2) This is also a solution, but $$f(Q)f^{}(Q)$$ (3) (cf. $`𝐠(x^{})𝐠^{}(x^{})`$).<sup>19</sup><sup>19</sup>19It would have been more natural to denote $`Q`$ by $`P^{}`$ and later to perform an active transformation dropping $`P^{}`$ (cf. $`x^{}`$) instead of $`P`$ (cf. $`x`$), in which case one would have obtained the expressions $`f(P)=f^{}(P^{})`$ instead of (1) and $`f(P)f^{}(P)`$ instead of (3), more in tune with Einstein’s original coordinate notation. However, remaining faithful to this starting point would have somewhat obscured the notation in the investigation that follows at the end of this section. Again, in order to preserve the causality of the field equation, one postulates diffeomorphism invariance, i.e. that “displacements” of points lead to no observable effects and that, therefore, points have no physical reality. How to move forward, then? We note that the meaning of the postulate stated above is only at face value so. On closer inspection, diffeomorphism invariance involves only the weaker requirement that points be all alike, i.e. have no physical identity; thereby, the problem posed by the hole argument is avoided, since having physically indistinguishable spacetime points makes $`𝐠(P)`$ and $`\varphi ^{}[𝐠(P)]`$ physically equivalent. In order to appreciate why having physically indistinguishable points solves the hole argument, one must clearly recognize the difference between mathematical and physical points. Mathematical points may very well be labelled points but, as Stachel remarked, > \[N\]o mathematical coordinate system is physically distinguished per se; and without such a distinction there is no justification for physically identifying the points of a \[mathematical\] manifold…as physical events in space-time. Thus, the mathematician will always correctly regard the original and the dragged-along fields as distinct from each other. But the physicist must examine this question in a different light…(Stachel, 1989, p. 75) The physicist must, in this case, rather ask whether there is anything in empty spacetime by means of which its points can be physically told apart from one another—neither labels nor metrics can count to carry out the differentiation. Finding there are no such means, the physicist must hold spacetime points physically identical to one another. However, having a multitude of physically indistinguishable spacetime points does not necessarily mean that they must be physically meaningless in every other way (cf. hydrogen atoms; see below). An earlier attempt by one of us (Lehto, Nielsen, & Ninomiya, 1986a,b) in which the principle of diffeomorphism invariance and quantum theory were both taken into account, revealed that fields on a pregeometric lattice displayed quantum-mechanical correlations. In that work, the quantum-mechanical framework was realized via a path-integral formalism; within it, the requirement of diffeomorphism invariance demanded that one ought to sum over all the vertices $`n`$ in the partition function $`Z`$, thus leading to the appearance of correlations. One the other hand, the said requirement induced a free gas behaviour<sup>20</sup><sup>20</sup>20Free gas behaviour in this case means that any pair of vertices can with high probability have any mutual distance. of the vertices, which helped to avoid the rise of long-range correlations.<sup>21</sup><sup>21</sup>21The correlation function $`\|f(P)f(Q)f(P)f(Q)\|`$ for fields on the lattice was found to be non-null, although correlations were nevertheless semi-local, i.e. actually tending to zero exponentially as the distance between two lattice points $`P`$ and $`Q`$ increased. Although we do not plan to follow this earlier approach, we rescue from it the insinuation that the said principle and quantum theory—the only branch of natural science so far forced to confront the problem of existence<sup>22</sup><sup>22</sup>22See (Isham, 1995, p. 65).—can lead together to the result that physical fields can be mutually correlated in spacetime. We hold such correlations to be the key to the possible identification of spacetime observables. An analogy with matter that behaves quantum-mechanically may clarify in what sense correlations can fulfill this task. Just like spacetime points are physically identical, so are the hydrogen atoms conforming a gas of this element. However, denying the reality of hydrogen atoms on these grounds does not seem at all reasonable (Horwich, 1978, p. 409; Friedman, 1983, p. 241). Hydrogen atoms, despite being identical, possess a property that spacetime points do not seem to have: they interact with other matter and correlate with one another creating bonds, hydrogen molecules. Holding the hydrogen molecules to be our sought-for observables for the sake of this analogy,<sup>23</sup><sup>23</sup>23In a genuinely empty spacetime, points could only “interact” among themselves. Therefore, the interaction of the hydrogen atoms with e.g. their container’s walls is an analogy we cannot pursue. their measurable properties now give evidence of the existence of the physically identical atoms. Forty years ago, Wheeler touched upon an idea somewhat similar to ours within the context of his so-called “bucket of dust.” He wrote: > Two points between which any journey was previously very long have suddenly found themselves very close together. However sudden the change is in classical theory, in quantum theory there is a probability amplitude function which falls off in the classically forbidden domain. In other words, there is some residual connection between points which are ostensibly very far apart. (Wheeler, 1964, p. 498) However, Wheeler did not develop this idea, but only used it in order to reject his quantum-mechanical concept of nearest neighbour, according to the manner in which he had previously defined it. In order to display the physically new idea of spacetime correlations within the existing context of the hole argument—but moving, at the same time, beyond it—we proceed as follows. In an otherwise empty spacetime, a field $`f(P)`$ at point $`P`$ dragged onto another point $`Q`$, $$\varphi ^{}[f(P)]=f[\varphi (P)]=f(Q),$$ (4) and then pulled back again, $`\varphi _{}[f(Q)]`$, could carry properties pertaining to $`Q`$ back with it, so that a comparison of $`f(P)`$ and $`\varphi _{}[f(Q)]`$ could yield that they are physically different. This is not to say that we expect to find new physics by means of a purely mathematical operation plus its inverse. It rather means that, given the insight that indistinguishable points may nonetheless display physical effects by correlating with each other quantum-mechanically, then the above analysis constitutes a means to represent the ensuing broken physical symmetry of diffeomorphism invariance. Moreover, when we speak of “moving points,” no physical system is actually being displaced in the physical world; the physical meaning of the expressions above (and below) ultimately falls back on spacetime correlations, the existence of which is here conjectured. Given $`f(P)`$, the appearance of the pull-back $`\varphi _{}[f(Q)]`$ rests on the need to compare locally the original field and the original field restored. This need arises from the fact that physical experiments are not performed globally but locally. An analogy close at hand is that of parallel transport in general relativity. Given two vector fields $`\stackrel{}{v}(P)`$ and $`\stackrel{}{v}(Q)`$ in curved spacetime, they can only be compared by computing $`\mathrm{\Delta }\stackrel{}{v}`$ at one point $`P`$ by parallel-transporting the latter field: $$\mathrm{\Delta }\stackrel{}{v}(P)=\stackrel{}{v}(Q)_{}\stackrel{}{v}(P).$$ (5) The analogy works best for the case of a small, closed loop with sides $`\mathrm{\Delta }a\stackrel{}{e}_1`$ and $`\mathrm{\Delta }b\stackrel{}{e}_2`$. In this case, the field $`\stackrel{}{v}(P)`$ is compared to itself after having travelled the loop and returned to its original position; $$\mathrm{\Delta }\stackrel{}{v}(P)=\mathrm{\Delta }a\mathrm{\Delta }b𝐑(,\stackrel{}{v}(P),\stackrel{}{e}_1,\stackrel{}{e}_2),$$ (6) where $`𝐑`$ is the Riemann tensor, quantifies how much the components—but not the size—of $`\stackrel{}{v}(P)`$ have been changed by the experience. Similarly, if instead of thinking of a field as a mere value, we visualize it as a vector $`f`$, and $`f(P)`$ as a component of $`f`$, then a new dimension to the hole argument is revealed: the local comparison of $`\varphi _{}[f(Q)]`$ and $`f(P)`$ quantifies how much the field component $`f(P)`$—but not the field’s value—changes due to the possible correlation of $`P`$ with $`Q`$; in short, $`f`$ as a vector would behave like $`\stackrel{}{v}`$. This change would stem from any physical reality of quantum-mechanical origin that the points may have; pictorially speaking, $`P`$ would behave as if it “remembered” having “interacted” with $`Q`$, remaining “entangled” with it. Now in the same manner that the parallel-transport procedure constitutes the mathematical representation of a physical property of spacetime—namely, its curvature—so does our description intend to represent a new kind of physical property—namely, spacetime correlations. Further, just as curvature reveals itself as the active effect of spacetime’s geometric structure and may well be viewed as a sort of lingering connection between the points on the loop, the correlations we envision would likewise unveil themselves as the active effect of a persistent connection between spacetime points. However, in this case, we must search for the source of this connection somewhere else, within a deeper layer of the nature of things (see below). The challenge now is to find observables which reveal this possible behaviour of field $`f(P)`$, since fields alone are not observable as such. In this connection, we note with Einstein (1952b, pp. 119, 121, 131; 1970, p. 71; 1982, p. 47) that the intervals $`\mathrm{d}s^2`$—and not the metric field—are the fundamental constituents of general relativity: the theory is essentially about an observable network of invariant intervals between events. Through the intervals, a featureless spacetime acquires geometric structure, which then can be characterized via the metric tensor (and the inner product), thus: $`\mathrm{d}s^2=\mathrm{d}\stackrel{}{s}\mathrm{d}\stackrel{}{s}`$ $`=`$ $`\left({\displaystyle \frac{\stackrel{}{s}}{x^\mu }}\|_P\mathrm{d}x^\mu \right)\left({\displaystyle \frac{\stackrel{}{s}}{x^\nu }}\|_P\mathrm{d}x^\nu \right)`$ (7) $`=`$ $`\left({\displaystyle \frac{\stackrel{}{s}}{x^\mu }}\|_P{\displaystyle \frac{\stackrel{}{s}}{x^\nu }}\|_P\right)\mathrm{d}x^\mu \mathrm{d}x^\nu `$ $`=`$ $`g_{\mu \nu }(P)\mathrm{d}x^\mu \mathrm{d}x^\nu .`$ In keeping with this crucial realization about $`\mathrm{d}s^2`$, the interval between points suggests itself as the main candidate for a spacetime observable. One could explore<sup>24</sup><sup>24</sup>24Here $`P^{}`$ and $`Q^{}`$ are points in the neighbourhoods of $`P`$ and $`Q`$ respectively, and are linked by a diffeomorphism $`\varphi (P^{})=Q^{}`$, in the same way that $`P`$ and $`Q`$ are. whether the interval $$\mathrm{d}s_{PP^{}}^2=g_{\mu \nu }(P)\mathrm{d}x^\mu \mathrm{d}x^\nu $$ (8) between two points<sup>25</sup><sup>25</sup>25Special care must be taken here not to confuse points with events. remains unchanged after $`g_{\mu \nu }(P)`$ is pushed forward, $`\varphi ^{}[g_{\mu \nu }(P)]`$, to get $`\mathrm{d}s_{QQ^{}}^2`$ $`=`$ $`\varphi ^{}[g_{\mu \nu }(P)]\mathrm{d}x^\mu \mathrm{d}x^\nu `$ (9) $`=`$ $`g_{\mu \nu }(Q)\mathrm{d}x^\mu \mathrm{d}x^\nu ,`$ and subsequently pulled back, $`\varphi _{}[g_{\mu \nu }(Q)]`$, to get for the original interval $$\varphi _{}[g_{\mu \nu }(Q)]\mathrm{d}x^\mu \mathrm{d}x^\nu .$$ (10) If spacetime points had any physical reality, the final expression (10) could not be equal to the initial one (8), and there would be some long-range correlations seen in the line element. This would not upset general relativity because, as said above, the envisioned correlations represent a new kind of effect rather than corrections to already-known physical magnitudes, and therefore do not challenge the existing predictions of this theory. Moreover, at a conceptual level, the diffeomorphism invariance principle and the hole argument’s conclusion are not upset either, since these belong in general relativity’s classical, geometric description of spacetime, which does not consider quantum theory at all and, therefore, must view quantum-mechanical correlations as an element foreign to its framework. In particular and furthering the previous analogy, we may say that the local comparison between $`\varphi _{}[f(Q)]`$ and $`f(P)`$ deals with changes in the components of the field, whereas the hole argument refers to (lack of) changes in the value of the fields. Thus, the classical, geometric theory of general relativity remains untouched from this perspective. Indeed, consistent with the general relativistic picture, it is not, strictly speaking, possible for the correlations we envision to be displayed by $`\mathrm{d}s_{PP^{}}^2`$ as a geometric notion, i.e. as a distance between points. We interpret the appearance of quantum-mechanical correlations as an indication that there must be something amiss with the current geometric description—which leads to the futile problem of the geometric ether—and as evidence of physical things beyond geometry.<sup>26</sup><sup>26</sup>26In fact, the consideration of quantum-mechanical ideas in themselves involves, from our point of view, the need for non-geometric physical things. In a future article, we will present quantum theory on the basis of measurement results $`a_i`$, and metageometric premeasurement and transition things, $`𝒫(a_i)`$ and $`𝒫(a_j\|a_i)`$, familiar to human experience. We will argue that, when viewed in this manner, the theory gets rid of some of the philosophical problems that plague it (e.g. the geometric state vector $`\|\psi `$ and its controversial ontology; cf. the geometric points $`P`$ and their controversial ontology). In other words, we understand what we presently describe geometrically as “correlations between spacetime points” to be in fact the effect of quantum-mechanical, metageometric things. In particular, we expect the geometric interval $`\mathrm{d}s_{PP^{}}^2`$ to result as a geometric remnant or trace of such things. To visualize our meaning, consider the following extension of the earlier parallel-transport analogy. Due to its homogeneity, a flat Euclidean spacetime is sterile as far as observables related to its physical geometry are concerned; a curved (pseudo-)Riemannian spacetime, on the other hand is not: the parallel transport of a field around a closed loop uncovers a previously hidden geometric observable, this spacetime’s curvature. A similar relationship holds now between a vacated general-relativistic spacetime and what we envision as a metageometric realm: whereas the former is completely homogenous concerning its elementary points (diffeomorphism symmetry) and, therefore, sterile with respect to empty spacetime observables, the latter metageometric realm may uncover currently hidden spacetime observables as residual effects of the quantum-mechanical connection between metageometric things. It is through these things beyond geometry that we aspire to achieve the earlier anticipated conceptual overthrow of the geometric ether. This is the sense, then, in which we intend the question that serves as the title of this article, and in this sense only that we intend to pursue an answer to it. In this light, one should regard points the way one would Wittgenstein’s (1922) ladder: as concepts to be discarded after one “has climbed out through them, on them, over them” ($`\mathrm{\S }`$ 6.54). ## References 1. Auyang, S. Y. (2001). Spacetime as a fundamental and inalienable structure of fields. Studies in History and Philosophy of Modern Physics, 32, 205–215. 2. Brans, C. H. (1999). Absolute spacetime: The twentieth century ether. General Relativity and Gravitation, 31, 597-607, arXiv:gr-qc/9801029. 3. Butterfield, J. (1989). The hole truth. The British Journal for the Philosophy of Science, 40, 1–28. 4. Drude, P. (1894). Physik des Äthers auf elektromagnetisher Grundlage. Stuttgart: Enke. 5. Earman, J. (1989). World enough and space-time: Absolute versus relational theories of space and time. 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Boston: Birkhäuser (Reprinted from Archive for History of Exact Sciences, 38 (1988)). 19. Lehto, M., Nielsen, H. B., & Ninomiya, M. (1986a). Pregeometric quantum lattice: A general discussion. Nuclear Physics B, 272, 213–227. 20. Lehto, M., Nielsen, H. B., & Ninomiya, M. (1986b). Diffeomorphism symmetry in simplicial quantum gravity. Nuclear Physics B, 272, 228–252. 21. Meschini, D., Lehto, M., & Piilonen, J. (2005). Geometry, pregeometry and beyond. Studies in History and Philosophy of Modern Physics, 36, 435–464, arXiv:gr-qc/0411053. 22. Newton, I. (1962). Mathematical principles of natural philosophy (A. Motte, & F. Cajori, Trans.). Berkeley: University of California Press (Original work published 1729). 23. Rynasiewicz, R. (1994). The lessons of the hole argument. The British Journal for the Philosophy of Science, 45, 407–436. 24. Rynasiewicz, R. (1996). Is there a syntactic solution to the hole problem? Philosophy of Science, 63 (Proceedings), S55–S62. 25. Smolin, L. (2004, January). Atoms of space and time. Scientific American, pp. 56–65. 26. Stachel, J. (1989). Einstein’s search for general covariance. In D. Howard, & J. Stachel (Eds.), Einstein and the history of general relativity (pp. 63–100). Boston: Birkhäuser (Original work published 1980). 27. Weyl, H. (1918). Raum, Zeit, Materie. Berlin: Springer. 28. Weyl, H. (1949). Philosophy of mathematics and natural science. Princeton, NJ: Princeton University Press. 29. Wheeler, J. A. (1964). Geometrodynamics and the issue of the final state. In C. De Witt, & B. S. De Witt (Eds.), Relativity, groups and topology (pp. 317–520). New York: Gordon and Breach. 30. Whittaker, E. T. (1951). A history of the theories of aether and electricity (Vols. 1–2). London, New York: Tomash Publishers & American Institute of Physics. 31. Wittgenstein, L. (1922). Tractatus logico-philosophicus (C. K. Ogden, Trans.). London: Routledge & Kegan.
warning/0506/hep-ph0506018.html	ar5iv	text	# Long-distance contributions in 𝐾→𝜋⁢𝑙⁺⁢𝑙⁻ Decays ## 1 Introduction Nowadays, rare K decays know a particularly increasing interest since the branching ratios of some still unobserved of them, that are golden modes for CP–violation study, are predicted to be rather closed to actual experimental limits. Two of the most known and most promising processes belonging to this class are $`K_L\pi ^0\nu \overline{\nu }`$ and $`K_L\pi ^0e^+e^{}`$. In $`^\mathrm{?}`$, one of the aims was the evaluation of the $`K_L\pi ^0e^+e^{}`$ long-distance contributions. Until 2003 and the first observation and precise measurement by the NA48 Collaboration at CERN $`^\mathrm{?}`$ of the modes $`K_S\pi ^0e^+e^{}`$ and $`K_L\pi ^0\gamma \gamma `$, no really reliable theoretical estimate of the $`K_L\pi ^0e^+e^{}`$ branching ratio was available, since this quantity was thought, at least from some viewpoint, to depend crucially on these two decays. Then, after these new measurements, a theoretical re-analysis of this golden mode was obtained in $`^\mathrm{?}`$. Our work $`^\mathrm{?}`$ propose an alternative strategy for part of the theoretical analysis elaborated in this last paper (and in the correlated paper $`^\mathrm{?}`$). We shall not dive here into theoretical details that can be found in our original publication (see also $`^\mathrm{?}`$) but instead insist on some of the underlying ideas. Apart from its CP–conserving part which is known to be negligible, the branching ratio of the $`K_L\pi ^0e^+e^{}`$ process fragments into three parts: a direct CP–violating (CPV) part, an indirect one (also called mixing) and the important interference between them, $`\mathrm{Br}\left(K_L\pi ^0e^+e^{}\right)\|_{\mathrm{CPV}}`$ $`=\left[C_{dir.}\left({\displaystyle \frac{\mathrm{Im}\lambda _t}{10^4}}\right)^2+C_{ind.}\pm C_{int.}{\displaystyle \frac{\mathrm{Im}\lambda _t}{10^4}}\right]\times 10^{12},`$ (1) where $`C_{dir.}=2.4\pm 0.2`$ is known precisely due to its (perturbative) short-distance origin, $`C_{ind.}=10^{12}\|ϵ\|^2\frac{\tau (K_L)}{\tau (K_S)}\mathrm{Br}(K_S\pi ^0e^+e^{})`$ and $`C_{int.}\sqrt{C_{dir.}C_{ind.}}`$ (see $`^\mathrm{?}`$). With NA48 results concerning $`\mathrm{Br}(K_S\pi ^0e^+e^{})`$, we can directly conclude that the indirect CPV contribution is quite large. This apparently overwhelming CP violating behaviour of $`K_L\pi ^0e^+e^{}`$ explains the importance of this decay mode and thus one major question concerns the type of interference (constructive or destructive, depending on the algebraic value of $`C_{int.}`$). Arguments in favour of a constructive interference have been suggested in $`^\mathrm{?}`$ and are good news for future experimental investigations; as we shall see in the following, we found the same conclusion within the approach that we suggested. ## 2 Indirect CPV and interference contributions: $`𝒪(p^4)`$ Chiral Perturbation theory and beyond In our work, two of the main points of study are the indirect CPV and interference contributions of $`\mathrm{Br}(K_L\pi ^0e^+e^{})`$. We wish to emphasize that contrary to $`^\mathrm{?}`$ who rely on the $`K_S\pi ^0e^+e^{}`$ channel for prediction of the indirect CPV term, we showed in $`^\mathrm{?}`$ that both contributions can be related to the much more experimentally known charged mode $`K^+\pi ^+e^+e^{}`$ for which our approach gives also a description consistent with all corresponding available data. It results from this viewpoint a semi-theoretical prediction for the indirect CPV contribution competitive to the purely experimental deduction of $`^\mathrm{?}`$. Moreover, our approach, by its common framework for the neutral and charged modes, has less free parameters than what the authors of $`^\mathrm{?}`$ had to introduce in order to obtain at the same time their conclusion for the interference term and a theoretical description of the charged mode compatible with experiments (see also $`^\mathrm{?}`$ for part of the original work concerning the latter). The low-energy description of $`K\pi \mathrm{}^+\mathrm{}^{}`$ decays is based on Chiral Perturbation Theory ($`\chi `$PT) and, for the specific case of $`K_L\pi ^0e^+e^{}`$, the CPV–branching ratio (1) can be rewritten, at $`𝒪(p^4)`$ of $`\chi `$PT, as <sup>a</sup><sup>a</sup>aThis equation is obtained by keeping only $`𝒪(p^4)`$ terms and neglecting constant and linear terms in $`a_s`$ in eq. (30) of $`^\mathrm{?}`$. $`\mathrm{Br}\left(K_L\pi ^0e^+e^{}\right)\|_{\mathrm{CPV}}`$ $`=[(2.4\pm 0.2)\left({\displaystyle \frac{\mathrm{Im}\lambda _t}{10^4}}\right)^2+(3.5\pm 0.1)({\displaystyle \frac{1}{3}}\text{w}_s)^2`$ $`+(2.9\pm 0.2)({\displaystyle \frac{1}{3}}\text{w}_s){\displaystyle \frac{\mathrm{Im}\lambda _t}{10^4}}]\times 10^{12}.`$ (2) As can be seen in this formula, the type of interference is directly correlated to one unknown constant $`𝐰_s`$ that we have to evaluate (the other constant $`\text{Im}\lambda _t`$ is known to be $`(1.36\pm 0.12)\times 10^4`$, see $`^\mathrm{?}`$). With (2) in mind, we can now basically explain the strategy that has been proposed in our work in order to obtain a prediction for this branching ratio but in this respect, we have to tell more about $`𝐰_s`$. In fact, $`𝐰_s`$ is mainly the combination of low-energy chiral coupling constants that appear in the renormalized $`𝒪(p^4)`$ $`\chi `$PT decay rate’s expression of the $`K_S\pi ^0e^+e^{}`$ neutral mode $`^{\mathrm{?},\mathrm{?}}`$, $$𝐰_s=\frac{1}{3}(4\pi )^2\stackrel{~}{𝐰}\frac{1}{3}\mathrm{log}\frac{M_K^2}{\nu ^2},$$ (3) where $`\stackrel{~}{𝐰}=𝐰_1𝐰_2`$. It is possible to have a completely theoretical estimate of these coupling constants by finding their underlying Green’s functions that can then be evaluated in the Large-$`N_c`$ expansion of QCD (see, for an example of such procedure, ref. $`^\mathrm{?}`$). We first followed another, more phenomenological, approach. One should notice that, up to the logarithm, the same coupling constants combination also appears in the $`𝒪(p^4)`$ $`\chi `$PT expression of the $`K^+\pi ^+e^+e^{}`$ decay rate, in the corresponding $`𝐰_+`$ constant $`^{\mathrm{?},\mathrm{?}}`$ $$𝐰_+=\frac{1}{3}(4\pi )^2\stackrel{~}{𝐰}(4\pi )^2[𝐰_24𝐋_9]\frac{1}{6}\mathrm{log}\frac{M_K^2m_\pi ^2}{\nu ^4}.$$ (4) This is one of the two important points that make that the prediction for the indirect CPV and interference terms of $`\mathrm{Br}(K_L\pi ^0e^+e^{})`$ can be obtained entirely from the charged mode (the second point is that our ”beyond $`𝒪(p^4)`$$`\chi `$PT approach does not introduce more free parameters than what ”pure $`𝒪(p^4)`$$`\chi `$PT does, so that our unique unknowns in the ”beyond $`𝒪(p^4)`$$`\chi `$PT approach are still $`\stackrel{~}{𝐰}`$ and $`𝐰_24𝐋_9`$). Now, this is precisely the charged mode which is at the origin of the completely different theoretical analyses done in $`^\mathrm{?}`$ and in our work. In fact, two experimental informations are available for the charged decay: its branching ratio and the mass spectrum $`^\mathrm{?}`$. Now, $`\chi `$PT at $`𝒪(p^4)`$ fails to correctly reproduce the slope of the mass spectrum. Two different strategies can then be followed to solve this problem. The first one $`^\mathrm{?}`$ (see also $`^\mathrm{?}`$) is to consider that $`𝒪(p^6)`$ corrections are important and that it is necessary to go to this order to have a good description of the form factor that enters into the theoretical expression of the $`K^+\pi ^+e^+e^{}`$ decay rate. This is what the authors of $`^\mathrm{?}`$ have done in an approximate way for the charged channel (and as a by product for the neutral case) since the $`𝒪(p^6)`$ $`\chi `$PT analysis is too complicated to be done without approximation. With their approach, they were able to fit the $`K^+\pi ^+e^+e^{}`$ mass spectrum and their analysis of the neutral mode has then been reinvestigated in $`^\mathrm{?}`$ under the light of the NA48 new measurements, from which it has been possible to obtain a prediction for the indirect CPV and interference parts of $`\mathrm{Br}(K_L\pi ^0e^+e^{})`$. The second strategy to solve the $`𝒪(p^4)`$ problem is the one that has been put forward in our work. We proposed in $`^\mathrm{?}`$, with the help of a Large-$`N_c`$ QCD inspired model, a second way of dealing with the $`𝒪(p^4)`$ problem by slighty modifying the $`𝒪(p^4)`$ form factor expression with an explicit replacement of the $`𝒪(p^4)`$ couplings by their underlying minimal resonances structure. This achieves a resummation to all orders of $`\chi `$PT which then goes beyond $`𝒪(p^6)`$. More precisely, we replaced the $`𝒪(p^4)`$ form factor $`^{\mathrm{?},\mathrm{?}}`$ $$f_V(z)=\frac{G_8}{G_F}\left\{\frac{1}{3}𝐰_+\frac{1}{60}z\chi (z)\right\},$$ (5) with $`\chi (z)=\varphi _\pi (z)\varphi _\pi (0)`$ and $`\varphi _\pi (z)=\frac{4}{3}\frac{m_\pi ^2}{M_K^2z}+\frac{5}{18}+\frac{1}{3}(\frac{4m_\pi ^2}{M_K^2z}1)^{\frac{3}{2}}\mathrm{arctan}(\frac{4m_\pi ^2}{M_K^2z}1)^{\frac{1}{2}}`$, by the formula $`f_V(z)`$ $`={\displaystyle \frac{G_8}{G_F}}\{{\displaystyle \frac{(4\pi )^2}{3}}[\stackrel{~}{𝐰}{\displaystyle \frac{M_\rho ^2}{M_\rho ^2M_K^2z}}+6F_\pi ^2\beta {\displaystyle \frac{M_\rho ^2M_K^{}^2}{\left(M_\rho ^2M_K^2z\right)\left(M_K^{}^2M_K^2z\right)}}]`$ $`+{\displaystyle \frac{1}{6}}\mathrm{ln}\left({\displaystyle \frac{M_K^2m_\pi ^2}{M_\rho ^4}}\right)+{\displaystyle \frac{1}{3}}{\displaystyle \frac{1}{60}}z\chi (z)\},`$ (6) where $`\beta =\frac{M_\rho ^2}{2F_\pi ^2}(1\frac{M_\rho ^2}{M_K^{}^2})^1(𝐰_24𝐋_9)`$. We kept the lowest order chiral loop contribution as the leading manifestation of the Goldstone dynamics. As one can notice, no new free parameter appears in (2) compared to (5). Now, with $`\stackrel{~}{𝐰}`$ and $`\beta `$ left as free parameters, we can make a non-linear regression to the mass spectrum data of $`^\mathrm{?}`$. This is a typical errors-in-variables statistical problem but, in a first approximation, we neglected the invariant mass error of the data and the particles masses errors appearing in (2), planning to take into account these effects in a future statistical study that will also include the not yet available but forthcoming NA48 mass spectrum data. The result is the continuous curve shown in the Figure 1, which corresponds to a $`\chi _{\text{min.}}^2=13.0`$ for 18 degrees of freedom. The fitted values of the parameters (using $`𝐠_8=3.3`$ and $`F_\pi =92.4\text{MeV}`$) are $$\stackrel{~}{𝐰}=0.045\pm 0.003\text{and}\beta =2.8\pm 0.1.$$ (7) We then compute the $`K^+\pi ^+e^+e^{}`$ branching ratio, using (2) and the fitted values for $`\stackrel{~}{𝐰}`$ and $`\beta `$, with the result $$\mathrm{Br}(K^+\pi ^+e^+e^{})=(3.0\pm 1.1)\times 10^7,$$ (8) in good agreement with experiment result $$\mathrm{Br}(K^+\pi ^+e^+e^{})=(2.88\pm 0.13)\times 10^7.$$ (9) A form factor similar to (2) can be obtained for the case of $`\mathrm{Br}(K_S\pi ^0e^+e^{})`$ and making use of (7), we find $`\mathrm{Br}(K_S\pi ^0e^+e^{})`$ $`=`$ $`(7.7\pm 1.0)\times 10^9`$ (10) and $`\mathrm{Br}(K_S\pi ^0e^+e^{})\|_{>165\text{MeV}}`$ $`=`$ $`(4.3\pm 0.6)\times 10^9.`$ (11) This is to be compared with the NA48 results $`^\mathrm{?}`$ $$\mathrm{Br}(K_S\pi ^0e^+e^{})=[5.8_{2.3}^{+2.8}(\mathrm{stat}.)\pm 0.8(\mathrm{syst}.)]\times 10^9,$$ (12) and $$\mathrm{Br}(K_S\pi ^0e^+e^{})\|_{>165\text{MeV}}=[3_{1.2}^{+1.5}(\mathrm{stat}.)\pm 0.1(\mathrm{syst}.)]\times 10^9.$$ (13) The predicted branching ratios for the $`K\pi \mu ^+\mu ^{}`$ modes are $$\mathrm{Br}(K^+\pi ^+\mu ^+\mu ^{})=(8.7\pm 2.8)\times 10^8\text{and}\mathrm{Br}(K_S\pi ^0\mu ^+\mu ^{})=(1.7\pm 0.2)\times 10^9,$$ (14) to be compared with $`\mathrm{Br}(K^+\pi ^+\mu ^+\mu ^{})`$ $`=`$ $`(7.6\pm 2.1)\times 10^8,\text{ref. }^\mathrm{?}`$ (15) and $`\mathrm{Br}(K_S\pi ^0\mu ^+\mu ^{})`$ $`=`$ $`[2.9_{1.2}^{+1.4}(\mathrm{stat}.)\pm 0.2(\mathrm{syst}.)]\times 10^9,\text{ref. }^\mathrm{?}.`$ (16) Finally, the resulting negative value $`𝐰_s=2.1\pm 0.2`$ in (3) implies a constructive interference in the ”pure $`𝒪(p^4)`$$`\chi `$PT expression (2) with a predicted branching ratio $$\mathrm{Br}(K_L\pi ^0e^+e^{})\|_{\mathrm{CPV}}=(3.4\pm 0.4)\times 10^{11},$$ (17) where we have used $`^\mathrm{?}`$ $`\text{Im}\lambda _t=(1.36\pm 0.12)\times 10^4`$. When taking into account the effect of the modulating form factor by using (1) and (10), one finds $$\mathrm{Br}(K_L\pi ^0e^+e^{})\|_{\mathrm{CPV}}=(3.8\pm 0.4)\times 10^{11},$$ (18) which indicates that higher order terms in the chiral expansion seem not too important. ## 3 Conclusions Earlier analyses of $`K\pi e^+e^{}`$ decays within the framework of $`\chi `$PT have been extended beyond the predictions of $`𝒪(p^4)`$, by replacing the local couplings which appear at that order by their underlying narrow resonance structure in the spirit of the MHA to Large-$`N_c`$ QCD. The resulting modification of the $`𝒪(p^4)`$ form factor is very simple and does not add new free parameters. Taking as input the invariant $`e^+e^{}`$ mass spectrum only, our strategy allows to reproduce, within errors, all known $`K\pi \mathrm{}^+\mathrm{}^{}`$ branching ratios. The predicted interference between the direct and indirect CP–violation amplitudes in $`K_L\pi ^0e^+e^{}`$ is constructive, with an expected branching ratio (see (18)) within reach of a dedicated experiment. ## 4 Acknowledgments We are grateful to Eduardo de Rafael who is at the origin of this work and for his remarks concerning the manuscript. We are also grateful to Jérôme Charles for his help in the beginning implementation of our future statistical analysis. One of the authors, D.G., expresses his gratitude to Antonio Pich for the invitation to the Moriond Conference. ## 5 References
warning/0506/math0506255.html	ar5iv	text	# Large-deviations/thermodynamic approach to percolation on the complete graph ## 1. Introduction For physical systems, mean-field theory often provides a qualitatively correct description of “realistic behavior.” The corresponding analysis usually begins with the derivation of so called mean-field equations which are self-consistent relations involving the physical quantity of primary interest and the various parameters of the model. This approach may be realized and, to some extent, justified mathematically by considering the model on the complete graph where each constituent interacts with all others. As an example, let us consider the Ising model on a complete graph $`K_n`$ of $`n`$ vertices. Here we have a collection of $`\pm 1`$-valued random variables $`(\sigma _i)_{i=1}^n`$ which are distributed according to the probability measure $`\mu _n(\{\sigma \})=\text{e}^{\beta H_n(\sigma )}/Z_{n,\beta }`$, where $$H_n(\sigma )=\frac{1}{n}\underset{i,j=1}{\overset{n}{}}\sigma _i\sigma _jh\underset{i=1}{\overset{n}{}}\sigma _i$$ (1.1) and where $`\beta ,h`$ are parameters. The relevant physical quantity is the empirical magnetization, $`m_n(\sigma )=n^1_{i=1}^n\sigma _i`$. In terms of this quantity, $`H_n(\sigma )=\frac{1}{2}n[m_n(\sigma )]^2hm_n(\sigma )`$ and so $$𝔼_n(\sigma _1\|\sigma _j:j1)=\mathrm{tanh}\left[\beta (m_n(\sigma )+h)\right]+O(1/n).$$ (1.2) This permits the following “cavity argument:” Supposing that $`m_n`$ tends, as $`n\mathrm{}`$, to a value $`m_{}`$ in probability, we have that $`m_{}=lim_n\mathrm{}𝔼_n(\sigma _1)`$ obeys $$m_{}=\mathrm{tanh}\left[\beta (m_{}+h)\right].$$ (1.3) This is the mean-field equation for the (empirical) magnetization. Of course, the concentration of the law of $`m_n`$ still needs to be justified; cf for details. In the context of percolation , the relevant mean-field model goes under the name the Erdös-Renyi Random Graph. Here each edge of $`K_n`$ is independently occupied with probability $`\alpha `$$`/`$$`n`$, where $`0\alpha <\mathrm{}`$, and vacant with probability $`1\alpha /n`$. The relevant “physical” quantity is the giant-component density $`\varrho _{}`$, i.e., the limiting fraction of the vertices that belong to the giant component of the graph. The corresponding mean-field equation, $$\varrho _{}=1\text{e}^{\alpha \varrho _{}},$$ (1.4) is also readily derived from heuristic “cavity” considerations. As is well known, $`\varrho _{}=0`$ is the only solution for $`\alpha \alpha _\text{c}=1`$, while for $`\alpha >\alpha _\text{c}`$ there is another, strictly positive solution. This solution tends to zero as $`\alpha \alpha _\text{c}`$; hence we may speak of a continuous transition. While (1.31.4) are indeed straightforward to derive, matters at the level of mean-field equations are not always satisfactory; the problem being the existence multiple solutions. As it turns out, for the percolation model (as well as the $`k`$-core percolation) the proper choice is always the maximal solution, but prescriptions of this sort generically fail, e.g., for the Ising model (1.3) with $`h<0`$ and, as often as not, whenever there is a first-order transition. Thus, one is in need of an additional principle which determines which of the solutions is relevant. The existing mathematical approach to these difficulties—e.g., for percolation , see also , or the $`k`$-core —is to work with sufficient precision until the mean-field conclusions are rigorously established. Another approach—which admits some prospects of extendability beyond the complete graph —is to supplement the picture by the introduction of the mean-field free-energy function. For the Ising model, this is a function $`m\mathrm{\Phi }_{\beta ,h}(m)`$ such that $$\mu _n\left(m_n(\sigma )m\right)=\text{e}^{n\mathrm{\Phi }_{\beta ,h}(m)+o(n)},n\mathrm{},$$ (1.5) i.e., $`m\mathrm{\Phi }_{\beta ,h}(m)`$ is the large-deviation rate function for the probability of observing the event $`\{m_n(\sigma )m\}`$. This spells the end of the story from the perspective of probability and/or theoretical physics: One seeks the minimum of the free energy function, setting its derivative to zero yields the mean-field equations with the irrelevant solutions corresponding to the local extrema which are not absolute minima; see again . The free-energy approach to mean-field problems has met with success in Ising systems and, to some extent, it has been applied to the Potts and random-cluster models . However, no attempt seems to have been made to extend this technology to “purely geometrical” problems on the complete graph, specifically, ordinary percolation or $`k`$-core percolation. The purpose of this note is to derive the large-deviation rate function for the event that the random graph contains a fraction $`\varrho `$ of vertices in “large” components. As we will see, the function has a unique minimum for all $`\alpha `$ which coincides with the “correct” solution of (1.4). We do not necessarily claim that the resultant justification of this equation is easier than which already exists in the literature. However, the picture presented here provides some additional insights into the model while the overall approach indeed admits the possibility of generalizations. ## 2. Main results Consider the set of vertices $`𝒱=\{1,\mathrm{},n\}`$ and let $`(\omega _{kl})_{1k<ln}`$ be a collection of i.i.d. random variables taking value one with probability $`p`$ and zero with probability $`1p`$. Let $`=(\omega )`$ be the (random) set $`\{(k,l):1k<ln,\omega _{kl}=1\}`$. In accord with the standard notation, cf , we will use $`𝒢(n,p)`$ to denote the undirected graph with vertices $`𝒱`$ and edges $``$. Of particular interest are the cases where $`p`$ decays to zero proportionally to $`1`$$`/`$$`n`$. Since these are the only problems we will consider, let us set, for once and all, $`p=\alpha /n`$ for some fixed $`\alpha [0,\mathrm{})`$. We will denote the requisite probability measure by $`P_{n,\alpha }`$. In order to state our main theorems, we need to introduce some notation. First, consider the standard entropy function $$S(\varrho )=\varrho \mathrm{log}\varrho +(1\varrho )\mathrm{log}(1\varrho )$$ (2.1) and let $$\pi _1(\alpha )=1\text{e}^\alpha .$$ (2.2) In addition, consider the function $$\mathrm{\Psi }(\alpha )=\left(\mathrm{log}\alpha \frac{1}{2}\left[\alpha \frac{1}{\alpha }\right]\right)0$$ (2.3) and note that $`\mathrm{\Psi }(\alpha )<0`$ if and only if $`\alpha >1`$. Finally, let us also define $$\begin{array}{c}\mathrm{\Phi }(\varrho ,\alpha )=S(\varrho )\varrho \mathrm{log}\pi _1(\alpha \varrho )\hfill \\ \hfill (1\varrho )\mathrm{log}\left[1\pi _1(\alpha \varrho )\right](1\varrho )\mathrm{\Psi }\left(\alpha (1\varrho )\right).\end{array}$$ (2.4) Then we have: ###### Theorem 2.1 Consider $`𝒢(n,\alpha /n)`$ and let $`𝒱_r`$ be the set of vertices that are in connected components of size larger than $`r`$. Then for every $`\varrho [0,1]`$, $$\underset{ϵ0}{lim}\underset{n\mathrm{}}{lim}P_{n,\alpha }\left(\|𝒱_{ϵn}\|=\varrho n\right)^{1/n}=\text{e}^{\mathrm{\Phi }(\varrho ,\alpha )}.$$ (2.5) An inspection of Lemma 6.2 reveals that, conditional on $`\{\|𝒱_{ϵn}\|=\varrho n\}`$, with $`ϵ>0`$, there will be only one “large” component with probability tending to one as $`n\mathrm{}`$. Fig. 1 shows the graph of $`\mathrm{\Phi }`$ for various values of $`\alpha `$ which is archetypal of free-energy functions in complete graph setting. The figure indicates a unique global minimum; direct, albeit arduous differentiation of (2.4) yields the fact that all local extrema satisfy the mean-field equation (1.4). The extremum at $`\varrho =0`$ is ruled out for $`\alpha >1`$ by noting that, under these conditions, the last term in (2.4) is strictly positive. The corresponding conclusion may also be extracted from the following probabilistic argument: Let $`m=\varrho n`$ and note that $`e^{nS(\varrho )}`$ is then the exponential growth-rate of $`\left(\genfrac{}{}{0pt}{}{n}{m}\right)`$. This allows us to write $$\text{e}^{n\mathrm{\Phi }(\varrho ,\alpha )}=e^{o(n)}\left(\genfrac{}{}{0pt}{}{n}{m}\right)\left[\pi _1(\alpha \varrho )\right]^m\left[1\pi _1(\alpha \varrho )\right]^{nm}\text{e}^{(nm)\mathrm{\Psi }(\alpha (1\varrho ))}.$$ (2.6) Neglecting the $`\mathrm{\Psi }`$-term (which provides a lower bound on $`\mathrm{\Phi }`$), one sees a quantity reminiscent of binomial distribution. Well known results on the latter inform us that the right-hand side is exponentially small unless $$\pi _1(\alpha \varrho )\frac{m}{n},$$ (2.7) i.e., unless $`\varrho `$ satisfies the mean-field equation (1.4). If $`\mathrm{\Psi }`$ is set to zero, there are degenerate minima for $`\alpha >1`$; however, the $`\mathrm{\Psi }`$-function will lift the degeneracy and, in fact, create a local maximum at $`\varrho =0`$ once $`\alpha >1`$. Meanwhile, in the region of the maximal solution, $`\mathrm{\Psi }`$ has vanished and the above mentioned approximation is exact. ###### Remarks 2.2 (1) A closely-related, but different problem to the one treated above has previously been studied using large-deviation techniques. Indeed, in , O’Connell derived the large-deviation rate function for the event that the largest connected component is of size about $`\kappa n`$. Note, however, that this does not restrict the total volume occupied by these component. For $`\kappa `$ close to $`\varrho _{}`$ from (1.4)—explicitly, as long as the complement of the large component has effective $`\alpha `$ less than $`1`$—O’Connell’s rate function coincides with ours. But once $`\kappa `$ is sufficiently small, his conditioning will lead to the creation of several large components whose total volume is such that their complement is effectively subcritical. Consequently, O’Connell never needs to address the central issue of our proof; namely, the decay rate of the probability that supercritical percolation has no giant components. (This is what gives rise to the term $`\mathrm{\Psi }`$ in (2.4) and the dashed portion of the graph in Fig. 1.) In fact, his rate function is basically a concatenation of many scaled copies of the undashed portion of the graph in Fig. 1. (2) While the $`\mathrm{\Psi }`$-term in (2.4) has a non-trivial effect on the large-deviation questions studied here, it does not play any role for events whose probability is of order unity (or is subexponential in $`n`$). This is because $`\mathrm{\Psi }`$ “kicks in” only for $`\varrho `$ away from the minimizing value. This is not the case for the $`k`$-core where the corresponding large-deviation analysis suggests that the analogous term “kicks in” right at the minimizer and may even affect the fluctuation scales. One way to bring $`\mathrm{\Psi }`$ out of the “realm of exponentially-improbable” for percolation would be to give each configuration a weight suppressing large components. However, we will not pursue these matters in the present note. (3) Our control of the rate function is not sharp enough to provide a detailed description of the critical region, i.e., the situations when $`\alpha =1+O(n^{1/3})`$. The corresponding analysis of the scaling phenomena inside the “critical window” has been performed in . On the other hand, for $`\alpha >1`$ one should be able to sharpen the control of the rate function near its minimum to derive a CLT for the fluctuations of the size of the giant component. Several ingredients enter our proof of Theorem 2.1 which are of independent interest. We state these as separate theorems. The first one concerns the exponential decay rate for the probability that the random-graph is (completely) connected: ###### Theorem 2.3 Let $`K`$ denote the event that $`𝒢(n,\alpha /n)`$ is connected. Then $$P_{n,\alpha }(K)=(1e^\alpha )^n\text{e}^{O(\mathrm{log}n)},n\mathrm{},$$ (2.8) where $`O(\mathrm{log}n)`$ is bounded by a constant times $`\mathrm{log}n`$ uniformly on compact sets of $`\alpha [0,\mathrm{})`$. We remark that Theorem 2.3 holds with $`\text{e}^{O(\mathrm{log}n)}`$ replaced by $`C(\alpha )+o(1)`$, see for a proof. However, the requisite steps seem far in excess of the derivation in Sect. 3. Furthermore, various pieces of Theorem 2.3 have been discovered, apparently multiple times, in ; cf also the discussion following Lemma 3.3. Next we present a result concerning the event that $`𝒢(n,\alpha /n)`$ contains no cycles. Such problems have been extensively studied under the conditions where this probability is $`O(1)`$, see e.g. . Our theorem concerns the large-deviation properties of this event: ###### Theorem 2.4 Let $`L`$ be the event that $`𝒢(n,\alpha /n)`$ contains no cycles. Then $$\underset{n\mathrm{}}{lim}P_{n,\alpha }(L)^{1/n}=\{\begin{array}{cc}\alpha \mathrm{exp}\left(\frac{\alpha }{2}+\frac{1}{2\alpha }\right),\hfill & \text{if }\alpha >1,\hfill \\ 1,\hfill & \text{otherwise}.\hfill \end{array}$$ (2.9) Strictly speaking, this result is not needed for the proof of our main theorem; it is actually used to derive the exponential decay for the probability of the event that $`𝒢(n,\alpha /n)`$ contains only “small” components. Surprisingly, the decay rates for these two events are exactly the same: ###### Theorem 2.5 Let $`L`$ be the event that $`𝒢(n,\alpha /n)`$ contains no cycles and let $`B_r`$ be the event that there are no components larger than $`r`$. Then $$\underset{r\mathrm{}}{lim}\underset{n\mathrm{}}{lim\; inf}P_{n,\alpha }(B_r)^{1/n}=\underset{ϵ0}{lim}\underset{n\mathrm{}}{lim\; sup}P_{n,\alpha }(B_{ϵn})^{1/n}=\underset{n\mathrm{}}{lim}P_{n,\alpha }(L)^{1/n}.$$ (2.10) Update: In the present paper we prove Theorem 2.4 using enumeration and generating-function techniques. Recently, a probabilistic approach has been developed by which we obtain an expansion of $`P_{n,\alpha }(L_n)`$ to quantities of order unity. One advantage of the new approach is that it also permits the analysis of the conditional measure $`P_{n,\alpha }(\|L_n)`$; see . To finish the discussion of our results, let us give some reason for the word “thermodynamic” in the title. The motivation comes from an analogy with droplet formation in systems at phase transition. Such situations have been studied extensively in the context of percolation and Ising (and Potts) model under the banner of “Wulff construction,” see for a review of these matters. One of the principal questions underlying Wulff construction is as follows: Compute the probability—and the characteristics of typical configurations carrying the event—that a given fraction of the system is in one thermodynamic state (e.g., liquid) while the rest is in another state (e.g., gas). It turns out that the typical configurations are such that the two phases separate; a droplet of one phase “floats” in the other phase. The requisite probability is then given by a large-deviation expression whose rate function is composed of three parts: the “surface” energy and entropy of the droplet, the rate function for the probability that the droplet is all in one phase, and the rate function for the probability that the complement of the droplet is in the other phase. In the case under study, the droplet is exactly the giant component and its weight is just the probability that all vertices in the droplet are connected to each other. The “surface” energy is (the log of) the probability that no vertex inside is connected to no vertex outside; the entropy is (the log of) the number of ways to choose the corresponding number of sites. The weight of the phase outside simply amounts to the probability that all remaining components are of submacroscopic scale. When the leading-order exponential decay rate of all of these contributions is extracted using Theorems 2.32.5, we get a quantity that only depends on the fraction of vertices taken by the droplet. The resulting expression is the one on the right-hand side of (2.6). ## 3. Everybody connected The goal of this section is to prove Theorem 2.3. Our proof is based on showing that the probability in (2.8) is exactly the same probability in a related, directed graph problem. For a collection of vertices $`𝒱_n=\{1,\mathrm{},n\}`$ and a set of edge probabilities $`(p_{kl})_{1k<ln}`$, let $`𝒢`$ be the inhomogeneous undirected random graph over $`𝒱_n`$. Similarly, let $`\stackrel{}{𝒢}`$ denote the inhomogeneous directed complete random graph with the restriction that the two possible (directed) edges between $`k`$ and $`l`$ occur independently, each with probability $`p_{kl}`$. To keep our notation distinct from the special case $`p_{kl}=\alpha /n`$ treated throughout this paper, we will write $`P`$ instead of $`P_{n,\alpha }`$. ###### Definition 3.1 A labelled directed graph $`𝒢=(𝒱,)`$ is said to be grounded at vertex $`v𝒱`$ if for every $`w𝒱`$ there exists a (directed) path from $`w`$ to $`v`$ in $``$. The identification of the two problems is now stated as follows: ###### Lemma 3.2 Let $`K`$ be the event that $`𝒢`$ is connected and let $`G`$ be the event that $`\stackrel{}{𝒢}`$ is grounded at vertex “$`1`$.” Then $`P(K)=P(G)`$. Proof. We use induction on the total number of edges incident with vertex “$`n`$.” Indeed, if $`p_{kn}=0`$ for all $`k=1,\mathrm{},n1`$, then $`P(K)=P(G)`$ because both probabilities are zero. Now let us suppose that $`P(K)=P(G)`$ when $`p_\mathrm{}n=0`$ for all $`\mathrm{}=k,\mathrm{},n1`$ and let us prove that it also for $`p_{kn}>0`$. It clearly suffices to show that the partial derivatives of $`P(K)`$ and $`P(G)`$ with respect to $`p_{kn}`$ are equal for all $`p_{kn}[0,1]`$. Notice first that both $`K`$ and $`G`$ are increasing events. Invoking Russo’s formula, see or \[22, Theorem 2.25\], we obtain $$\frac{}{p_{kn}}P(G)=P\left((n,k)\text{ is pivotal for }G\right),$$ (3.1) where the event $`\{(n,k)\text{ is pivotal for }G\}`$ means that if $`(n,k)`$ is occupied, the event $`G`$ occurs and if not, it does not. (Note that $`(n,k)`$ denotes the edge going from “$`n`$” to “$`k`$.”) The conditions under which this event occurs are straightforward: The set $`𝒱_n=\{1,\mathrm{},n\}`$ splits into two disjoint components, one rooted at “$`1`$” and the other at “$`n`$,” such that no vertex in the component associated with vertex “$`n`$” has an oriented edge to the other component and $`k`$ has an oriented path to $`1`$. Similarly, we have $$\frac{}{p_{kn}}P(K)=P\left((n,k)\text{ is pivotal for }K\right).$$ (3.2) Here $`\{(n,k)\text{ is pivotal for }K\}`$ simply means that, if the edge $`(n,k)`$ is absent, $`𝒱_n`$ consist of two connected components, one containing “$`1`$” and the other containing “$`n`$.” To see the equality of partial derivatives, we split both “pivotal” events according to the component containing the vertex “$`n`$.” If $`𝒲`$ is a set of vertices such that $`n𝒲`$ and $`1𝒲`$, let $`𝒢_{n,𝒲}`$ and $`𝒢_{1,𝒲}`$ be the restrictions of $`𝒢`$ to $`𝒲`$, and $`𝒱_n𝒲`$, respectively. Similarly, let $`\stackrel{}{𝒢}_{n,𝒲}`$ and $`\stackrel{}{𝒢}_{1,𝒲}`$ be the corresponding “components” of the oriented graph. Let $`K_{n,𝒲}`$ and $`K_{1,𝒲}`$ be the events that $`𝒢_{n,𝒲}`$ and $`𝒢_{1,𝒲}`$ are connected and let $`G_{n,𝒲}`$ and $`G_{1,𝒲}`$ be the events that $`\stackrel{}{𝒢}_{n,𝒲}`$ is grounded at “$`n`$” and that $`\stackrel{}{𝒢}_{1,𝒲}`$ is grounded at “$`1`$,” respectively. Since these pairs of events are independent, we have $$P\left((n,k)\text{ is pivotal for }G\right)=\underset{\begin{array}{c}𝒲:n𝒲\\ 1,k𝒲\end{array}}{}P(G_{1,𝒲})P(G_{n,𝒲})P(C_𝒲)\|_{p_{kn}=0},$$ (3.3) where $`C_𝒲`$ is the event that no vertex in $`𝒲`$ has a (directed) edge to $`𝒱_n𝒲`$. But the induction assumption tells us that $`P(G_{1,𝒲})=P(K_{1,𝒲})`$ and $`P(G_{n,𝒲})=P(K_{n,𝒲})`$, and the symmetry of edge probabilities for the directed graph tells us that $`P(C_𝒲)`$ is the probability that $`𝒢_{n,𝒲}`$, and $`𝒢_{1,𝒲}`$ are not connected by an edge in $`𝒢`$. Substituting these into (3.3), we get the right-hand side of (3.2). This completes the induction step. ∎ From now on, let $`K`$ and $`G`$ pertain to the specific random graphs $`𝒢(n,\alpha /n)`$ and $`\stackrel{}{𝒢}(n,\alpha /n)`$. We begin with upper and lower bounds on $`P_{n,\alpha }(K)`$: ###### Lemma 3.3 $`P_{n,\alpha }(K)\left(1\left(1\alpha /n\right)^{n1}\right)^{n1}`$. Proof. Let $`E`$ be the event—concerning the graph $`\stackrel{}{𝒢}(n,\alpha /n)`$—that every vertex except number “$`1`$” has at least one outgoing edge. Then $`GE`$ and so $$P_{n,\alpha }(G)P_{n,\alpha }(E)=\left(1\left(1\alpha /n\right)^{n1}\right)^{n1}.$$ (3.4) Invoking Lemma 3.2, this proves the desired upper bound. ∎ We remark that the upper bound in Lemma 3.3 has been discovered (and rediscovered) several times in the past. It seems to have appeared in for the first time and later in and also . A generalization to arbitrary connected graphs has been achieved in . ###### Lemma 3.4 $`P_{n,\alpha }(K)\left(1\left(1\alpha /n\right)^{n1}\right)^{n1}\frac{1}{n}`$. Proof. Consider the following events for directed random graph $`\stackrel{}{𝒢}(n,\alpha /n)`$: Let $`E`$ be the event that every vertex, except vertex number “$`1`$,” has at least one outgoing edge and let $`F`$ be the event every such vertex has exactly one outgoing edge. Since $`GE`$, we have $$P_{n,\alpha }(G)=P_{n,\alpha }(E)P_{n,\alpha }(G\|E).$$ (3.5) We claim that $$P_{n,\alpha }(G\|E)P_{n,\alpha }(G\|F).$$ (3.6) Indeed, let us pick an outgoing edge for each vertex different from “$`1`$,” uniformly out of all edges going out of that vertex, and let us color these edges red. Let $`G^{}`$ be the event that $`G`$ occurs using only the red edges. The distribution of red edges conditional on $`E`$ is the same as conditional on $`F`$. Hence $`P_{n,\alpha }(G\|E)P_{n,\alpha }(G^{}\|E)=P_{n,\alpha }(G^{}\|F)`$. But, on $`F`$, every available edge is red and so $`P_{n,\alpha }(G^{}\|F)=P_{n,\alpha }(G\|F)`$. Combining these inequalities, (3.6) is proved. The number of configurations that $`\stackrel{}{𝒢}(n,\alpha /n)`$ can take on $`F`$ is exactly $`(n1)^{n1}`$. On the other hand, the number of configurations which result in $`\stackrel{}{𝒢}(n,\alpha /n)`$ being grounded is $`a_n=n^{n2}`$—the number of labelled trees with $`n`$ vertices. Hence $$P_{n,\alpha }(G\|F)\frac{n^{n2}}{(n1)^{n1}}\frac{1}{n}.$$ (3.7) Using that $`P_{n,\alpha }(E)=(1(1\alpha /n)^{n1})^{n1}`$ the desired bound follows. ∎ Proof of Theorem 2.3. The claim is proved by noting $$\underset{n\mathrm{}}{lim}\frac{\left(1\left(1\alpha /n\right)^{n1}\right)^{n1}}{\left(1e^\alpha \right)^{n1}}=\mathrm{exp}\left((1\alpha /2)\frac{\alpha e^\alpha }{1e^\alpha }\right)$$ (3.8) and using the results of Lemmas 3.3 and 3.4. ∎ ## 4. Only trees Here we will assemble the necessary ingredients for the proof of Theorem 2.4. The proof is based on somewhat detailed combinatorial estimates and arguments using generating functions. Recall that $`L`$ denotes the event that $`𝒢(n,\alpha /n)`$ contains no cycles and that $`B_r`$ denotes the event that all components of $`𝒢(n,\alpha /n)`$ have no more than $`r`$ vertices. We begin by a combinatorial representation of the probability $`P_{n,\alpha }(LB_r)`$: Let $`a_{\mathrm{}}`$ denote the number of labeled trees on $`\mathrm{}`$ vertices. Then $`P_{n,\alpha }(LB_r)`$ $`={\displaystyle \underset{\begin{array}{c}{\scriptscriptstyle m_{\mathrm{}}\mathrm{}}=n\\ m_{\mathrm{}}=0\mathrm{}>r\end{array}}{}}{\displaystyle \frac{n!}{_{\mathrm{}}\left[m_{\mathrm{}}!(\mathrm{}!)^m_{\mathrm{}}\right]}}\left({\displaystyle \underset{\mathrm{}1}{}}\left[a_{\mathrm{}}\left({\displaystyle \frac{\alpha }{n}}\right)^\mathrm{}1\right]^m_{\mathrm{}}\right)\left(1{\displaystyle \frac{\alpha }{n}}\right)^{\left(\genfrac{}{}{0pt}{}{n}{2}\right)n+{\scriptscriptstyle m_{\mathrm{}}}}`$ (4.1) $`=n!\left({\displaystyle \frac{\alpha }{n}}\right)^n\left(1{\displaystyle \frac{\alpha }{n}}\right)^{\left(\genfrac{}{}{0pt}{}{n}{2}\right)n}{\displaystyle \underset{k=1}{\overset{n}{}}}\left({\displaystyle \frac{\alpha }{n}}\right)^k\left(1{\displaystyle \frac{\alpha }{n}}\right)^kQ_{n,k,r},`$ where we set $`k=_{\mathrm{}}m_{\mathrm{}}`$, applied the constraint $`_{\mathrm{}}\mathrm{}m_{\mathrm{}}=n`$ and let $`Q_{n,k,r}`$ denote the sum $$Q_{n,k,r}=\underset{\begin{array}{c}{\scriptscriptstyle m_{\mathrm{}}\mathrm{}}=n\\ {\scriptscriptstyle m_{\mathrm{}}}=k\\ m_{\mathrm{}}=0\mathrm{}>r\end{array}}{}\underset{\mathrm{}1}{}\left(\frac{a_{\mathrm{}}}{\mathrm{}!}\right)^m_{\mathrm{}}\frac{1}{m_{\mathrm{}}!}.$$ (4.2) We begin by isolating the large-$`n,k`$ behavior of this quantity: ###### Proposition 4.1 Consider the polynomial $$F_r(s)=\underset{\mathrm{}=1}{\overset{r}{}}\frac{s^{\mathrm{}}a_{\mathrm{}}}{\mathrm{}!}$$ (4.3) Then for all $`n,k,r1`$, $$Q_{n,k,r}\frac{1}{k!}\underset{s>0}{inf}\frac{F_r(s)^k}{s^n}.$$ (4.4) Moreover, for each $`\eta >0`$, there is $`n_0<\mathrm{}`$ and a sequence $`(c_r)_{r1}`$ of positive numbers for which $$Q_{n,k,r}\frac{c_r}{\sqrt{n}}\frac{1}{k!}\underset{s>0}{inf}\frac{F_r(s)^k}{s^n}$$ (4.5) holds for all $`nn_0`$, all $`k1`$ and all $`r2`$ such that $`k<(1\eta )n`$ and $`rk>n(1+\eta )`$. Proof of upper bound. Let us consider the generating function $$\widehat{Q}_r(s,z)=1+\underset{n=1}{\overset{\mathrm{}}{}}\underset{k=1}{\overset{n}{}}Q_{n,k,r}z^ks^n=\mathrm{exp}\left\{zF_r(s)\right\},$$ (4.6) where we used Fubini-Tonelli to derive the second equality. Since $`F_r`$ is a polynomial, the Cauchy integral formula yields $$Q_{n,k,r}=\frac{1}{(2\pi \text{i})^2}\text{d}s\text{d}z\frac{\mathrm{exp}\{zF_r(s)\}}{s^{n+1}z^{k+1}}=\frac{1}{2\pi \text{i}}\frac{1}{k!}\text{d}s\frac{F_r(s)^k}{s^{n+1}},$$ (4.7) where all integrals are over a circle of positive radius centered at the origin of $``$. Since all coefficients of $`F_r`$ are non-negative, $`\theta \|F_r(s\text{e}^{\text{i}\theta })\|`$ for $`s>0`$ is maximized at $`\theta =0`$. Bounding the integrand by its value at $`\theta =0`$, the integral yields a factor $`2\pi `$; optimizing over $`s>0`$ then gives the upper bound in (4.4). ∎ Proof of lower bound. As is common in Tauberian arguments, the lower bound will require somewhat more effort. First let us note that under the conditions $`k<(1\eta )n`$ and $`rk>n(1+\eta )`$ the function $`sF_r(s)^k/s^n`$, for $`s>0`$, blows up both at $`0`$ and $`\mathrm{}`$. Its minimum is thus achieved at an interior point; for the rest of this proof we will fix $`s`$ to a minimizer of this function. Since $`\|F_r(s\text{e}^{\text{i}\theta })\|<F_r(s)`$ for all $`\theta (\pi ,\pi ]\{0\}`$, the part of the integral in (4.7) corresponding to $`\|\theta \|>ϵ`$ is exponentially small (in $`n`$) compared to the infimum in (4.5). We thus need to show the lower bound only for the portion of the integral over $`\theta `$ with $`\|\theta \|ϵ`$, for some fixed $`ϵ>0`$. Since $`F_r`$ has positive coefficients, $`F_r0`$ in the (complex) $`ϵ`$-neighborhood of $`s`$. This allows us to define the function $$g(\theta )=\mathrm{log}\frac{F_r(s\text{e}^{\text{i}\theta })^\varrho }{s\text{e}^{\text{i}\theta }},\|\theta \|ϵ,$$ (4.8) where $`\varrho `$ plays the role of $`k`$$`/`$$`n`$. The function $`g`$ is analytic in an $`O(ϵ)`$-neighborhood of the origin. The choice of $`s`$ implies that $`g^{}(0)=0`$ which is equivalent to $$\frac{sF_r^{}(s)}{F_r(s)}=\frac{1}{\varrho }.$$ (4.9) For the second derivative we get $`g^{\prime \prime }(0)=\varrho \text{Var}(X)`$, where $`X`$ is the random variable with law $$P(X=\mathrm{})=\frac{1}{F_r(s)}\frac{a_{\mathrm{}}s^{\mathrm{}}}{\mathrm{}!},\mathrm{}=1,\mathrm{},r.$$ (4.10) In particular, since our restrictions on $`\varrho `$ between $`\frac{1}{r}(1+\eta )`$ and $`1\eta `$ imply that $`s`$ is bounded away from zero, this law is non-degenrate and so $`g^{\prime \prime }(0)<0`$. The analyticity of $`\theta g(\theta )`$ for $`\theta =O(ϵ)`$ implies that $`g^{\prime \prime \prime }`$ is bounded in this neighborhood, and so by Taylor’s theorem we have $$g(\theta )=g(0)A\theta ^2+O(\theta ^3),$$ (4.11) where $`A=A(r,\varrho )`$ is positive uniformly in the allowed range of $`\varrho `$’s and $`O(\theta ^3)`$ is a quantity bounded by $`\|\theta \|^3`$ times a constant depending only on $`r`$, $`ϵ`$ and $`\eta `$. (In particular, we may assume that $`O(\theta ^3)`$ is dominated by $`\frac{1}{2}A\theta ^2`$ for $`\|\theta \|ϵ`$.) We will split the integral over $`\theta [ϵ,ϵ]`$ into two more parts. Let $`\delta >0`$ and note that $`ng(0)`$ is the logarithm of the infimum in (4.5). Then for $`\theta `$ with $`\delta n^{1/3}\|\theta \|ϵ`$ we have $$n\text{Re}g(\theta )ng(0)\frac{1}{2}A\delta ^2n^{1/3}$$ (4.12) which shows that even this portion of the integral brings a contribution that is negligible compared to the right-hand side of (4.5). But for $`\|\theta \|\delta n^{1/3}`$ we have $`nO(\theta ^3)=O(\delta )`$ and so for $`\delta 1`$, the Taylor remainder will always have imaginary part between, say, $`\pi /4`$ and $`\pi `$$`/`$$`4`$. This means that $$\text{Re}_{\delta n^{1/3}}^{\delta n^{1/3}}\text{e}^{ng(\theta )}\text{d}\theta \frac{1}{2}\text{e}^{ng(0)}_{\delta n^{1/3}}^{\delta n^{1/3}}\text{e}^{nA\theta ^2}\text{d}\theta \frac{c}{\sqrt{n}}\text{e}^{ng(0)}$$ (4.13) for some constant $`c>0`$ which may depend on $`r`$ and $`\eta `$ but not on $`\varrho `$ and $`n`$. Combined with the previous estimates, this proves the lower bound (4.5). ∎ In light of the above lemma, the $`k`$-th term in the sum on the extreme right of (4.1) becomes $$\alpha ^kn^k\text{e}^{\alpha \frac{k}{n}}Q_{n,k,r}=\text{e}^{o(n)}\underset{s>0}{inf}\mathrm{exp}\left\{n\mathrm{\Theta }_r(s,k/n)\right\},$$ (4.14) where $$\mathrm{\Theta }_r(s,\varrho )=\varrho \mathrm{log}\alpha \varrho \mathrm{log}\varrho +\varrho +\varrho \mathrm{log}F_r(s)\mathrm{log}s.$$ (4.15) Here we should interpret (4.14) as an upper bound for $`r=n`$ and a lower bound for fixed $`r`$. It is clear that, regardless of $`r`$, the sum is dominated by $`k=\varrho n`$ for which $`\varrho inf_{s>0}\mathrm{\Theta }_r(s,\varrho )`$ is maximal. Such values are characterized as follows: ###### Lemma 4.2 Let $`\alpha >0`$ and $`r2`$. Then there is a unique $`(s_r,\varrho _r)[0,\mathrm{}]\times [1/r,1]`$ for which $$\mathrm{\Theta }_r(s_r,\varrho _r)=\underset{1/r\varrho 1}{sup}\underset{s>0}{inf}\mathrm{\Theta }_r(s,\varrho ).$$ (4.16) Moreover, we always have $`s_r(0,\mathrm{})`$ and $`\varrho _r(1/r,1)`$ and, furthermore, $$\underset{r\mathrm{}}{lim}\mathrm{\Theta }_r(s_r,\varrho _r)=\{\begin{array}{cc}1+\alpha /2\mathrm{log}\alpha ,\hfill & \text{if }\alpha 1,\hfill \\ 1+\frac{1}{2\alpha },\hfill & \text{if }\alpha >1.\hfill \end{array}$$ (4.17) Proof. We begin by ruling out the “boundary values” of $`s`$ and $`\varrho `$. First, if $`\varrho =1/r`$, then the infimum over $`s`$ is actually achieved by $`s=\mathrm{}`$. In that case $`F_r(s)=\mathrm{}`$ and the (one-sided) derivative with respect to $`\varrho `$ is infinite, i.e., $`\varrho =1/r`$ is a strict local minimum of $`\varrho inf_{s>0}\mathrm{\Theta }_r(s,\varrho )`$. Similarly, for $`\varrho =1`$ the infimum over $`s>0`$ is achieved at $`s=0`$ but then the $`\varrho `$-derivative of $`\varrho inf_{s>0}\mathrm{\Theta }_r(s,\varrho )`$ is negative infinity, i.e., also $`\varrho =1`$ is a strict local minimum. It follows that any $`(s_r,\varrho _r)`$ satisfying (4.16) necessarily lies in $`(0,\mathrm{})\times (1/r,1)`$. Setting the partial derivatives with respect to $`s`$ and $`\varrho `$ to zero shows that any minimizing pair is the solution of the equations $$F_r(s)=\alpha \varrho \text{and}sF_r^{}(s)=\alpha .$$ (4.18) In light of monotonicity of $`ssF_r^{}(s)`$, the solution is actually unique. To figure out the asymptotic as $`r\mathrm{}`$, we note that for $`s1/\text{e}`$, $$sF_r^{}(s)=\underset{\mathrm{}=1}{\overset{r}{}}a_{\mathrm{}}\frac{s^{\mathrm{}}}{(\mathrm{}1)!}\underset{r\mathrm{}}{}W(s),$$ (4.19) where $`W`$ is the unique number in $`[0,1/\text{e}]`$ such that $`W\text{e}^W=s`$. (Incidentally, $`W`$ is closely related to the survival probability of the Galton-Watson branching process with Poisson offspring distribution.) If $`s>1/\text{e}`$, then $`sF_r^{}(s)\mathrm{}`$ as $`r\mathrm{}`$. Using the relation between $`sF_r^{}(s)`$ and $`\alpha `$, we thus get $$s_r\underset{r\mathrm{}}{}\{\begin{array}{cc}\alpha \text{e}^\alpha ,\hfill & \text{if }\alpha 1,\hfill \\ 1/\text{e},\hfill & \text{if }\alpha >1.\hfill \end{array}$$ (4.20) Integrating the derivative of $`F_r`$ now shows that $`F_r(s_r)\alpha (1\alpha /2)`$ for $`\alpha 1`$. Using that $`F_r^{}(s)`$ is bounded for $`ss_r`$, we also find that $`F_r(s_r)1/2`$ for $`\alpha 1`$. This yields $$\varrho _r\underset{r\mathrm{}}{}\{\begin{array}{cc}1\alpha /2,\hfill & \text{if }\alpha 1,\hfill \\ \frac{1}{2\alpha },\hfill & \text{if }\alpha >1.\hfill \end{array}$$ (4.21) Noting that $`\mathrm{\Theta }(s_r,\varrho _r)=\varrho _r\mathrm{log}s_r`$ we now get (4.17). ∎ Proof of Theorem 2.4. By the fact that the supremum over $`\varrho `$ in (4.16) is achieved at an interior point, we can control the difference between the maximizing $`k`$$`/`$$`n`$ and its continuous counterpart $`\varrho `$. Thence $$P_{n,\alpha }(LB_r)=q_{n,r}n!\left(\frac{\alpha }{n}\right)^n\text{e}^{\alpha n/2}\mathrm{exp}\left\{n\mathrm{\Theta }_r(s_r,\varrho _r)\right\},$$ (4.22) where $$\frac{\stackrel{~}{c}_r}{\sqrt{n}}q_{n,r}n$$ (4.23) for some positive constants $`\stackrel{~}{c}_r`$ which may depend on $`r`$ and $`\alpha `$. Since $`B_n`$ contains every realization of $`𝒢(n,\alpha /n)`$, taking $`r=n`$ and applying Lemma 4.2 directly shows that $`P_{\alpha ,n}(L)\text{e}^{n\mathrm{\Psi }(\alpha )+o(n)}`$. To get a corresponding lower bound, we fix $`r2`$ and apply $`P_{\alpha ,n}(L)P_{\alpha ,n}(LB_r)`$. Taking $`1`$$`/`$$`n`$-th power and letting $`n\mathrm{}`$ then yields $$\underset{n\mathrm{}}{lim}P_{n,\alpha }(LB_r)^{1/n}=\alpha \text{e}^{1\alpha /2+\mathrm{\Theta }_r(s_r,\varrho _r)}.$$ (4.24) As we have just checked, the right-hand side tends to $`\text{e}^{\mathrm{\Psi }(\alpha )}`$ as $`r\mathrm{}`$. ∎ ###### Corollary 4.3 We have $$\underset{r\mathrm{}}{lim}\underset{n\mathrm{}}{lim}P_{\alpha ,n}(B_rL)^{1/n}=\underset{n\mathrm{}}{lim}P_{\alpha ,n}(L)^{1/n}.$$ (4.25) Proof. This summarizes the last step of the previous proof. ∎ ## 5. No big = no cycles Here we will prove that absence of large component has a comparable cost to absence of cycles, at least on an exponential scale. To achieve this goal, apart from Corollary 4.3, we will need the following upper bound: ###### Lemma 5.1 Let $`B_r`$ be the event that $`𝒢(n,\alpha /n)`$ has no components larger than $`r`$ and let $`L`$ be the event that all connected components of $`𝒢(n,\alpha /n)`$ are trees. Then for all $`r1`$, $$P_{n,\alpha }(B_r)P_{n,\alpha }(L)\left(1\frac{\alpha }{n}\right)^{\frac{1}{2}rn}.$$ (5.1) Proof. Let $`C`$ be the restriction of $`𝒢(n,\alpha /n)`$ to a set $`S\{1,\mathrm{},n\}`$. Let $`T`$ be a tree on $`S`$. Then $$\frac{P_{n,\alpha }(C=T)}{P_{n,\alpha }(CT)}=\left(1\frac{\alpha }{n}\right)^{\left(\genfrac{}{}{0pt}{}{\|S\|}{2}\right)\|S\|+1}\left(1\frac{\alpha }{n}\right)^{\frac{1}{2}\|S\|^2}.$$ (5.2) Hence $$P_{n,\alpha }(C\text{ is connected})\underset{T}{}P_{n,\alpha }(CT)\left(1\frac{\alpha }{n}\right)^{\frac{1}{2}\|S\|^2}P_{n,\alpha }(C\text{ is a tree}).$$ (5.3) Now, if $`L_r`$ is the event that no component of $`𝒢(n,\alpha /n)`$ of size larger than $`r`$ has cycles, then $`B_rL_r`$ and so $`P_{n,\alpha }(B_r)P_{n,\alpha }(L_r)`$. Let $`\{S_j\}`$ be a partition of $`\{1,\mathrm{},n\}`$ and let $`P_{n,\alpha }(\{S_j\})`$ denote the probability that $`\{S_j\}`$ are the connected components of $`𝒢(n,\alpha /n)`$. Then $$P_{n,\alpha }(L_r)=\underset{\{S_j\}}{}P_{n,\alpha }(\{S_j\})P_{n,\alpha }(L_r\|\{S_j\}),$$ (5.4) where $`P_{n,\alpha }(L_r\|\{S_j\})`$ is the conditional probability of $`L_r`$ given that $`\{S_j\}`$ are the connected components of $`𝒢(n,\alpha /n)`$. Letting $`C_j`$ represent the restriction of $`𝒢(n,\alpha /n)`$ to $`S_j`$, the bound (5.3) tells us that $`P_{n,\alpha }(L_r\|\{S_j\})`$ $`={\displaystyle \underset{j:\|S_j\|r}{}}P_{n,\alpha }(C_j\text{ is a tree}\|C_j\text{ is connected})`$ (5.5) $`{\displaystyle \underset{j}{}}P_{n,\alpha }(C_j\text{ is a tree}\|C_j\text{ is connected}){\displaystyle \underset{j:\|S_j\|<r}{}}\left(1{\displaystyle \frac{\alpha }{n}}\right)^{\frac{1}{2}\|S_j\|^2}`$ Using that $`\|S_j\|<r`$ for every $`S_j`$ contributing to the second product and applying that the sum of $`\|S_j\|`$ over the components with $`\|S_j\|<r`$ gives at most $`n`$, we then get $$P_{n,\alpha }(L_r\|\{S_j\})P_{n,\alpha }(L\|\{S_j\})\left(1\frac{\alpha }{n}\right)^{\frac{1}{2}rn}.$$ (5.6) Plugging this back in (5.4), the desired bound follows. ∎ Proof of Theorem 2.5. By Lemma 5.1 we have $$\underset{n\mathrm{}}{lim\; sup}P_{n,\alpha }(B_{ϵn})^{1/n}\text{e}^{ϵ/2}\underset{n\mathrm{}}{lim}P_{n,\alpha }(L)^{1/n}.$$ (5.7) On the other hand, the inclusion $`B_rB_rL`$ and Corollary 4.3 yield $$\underset{n\mathrm{}}{lim\; inf}P_{n,\alpha }(B_r)^{1/n}\underset{n\mathrm{}}{lim\; inf}P_{n,\alpha }(B_rL)^{1/n}\underset{r\mathrm{}}{}\underset{n\mathrm{}}{lim}P_{n,\alpha }(L).$$ (5.8) Since $`P_{n,\alpha }(B_r)P_{n,\alpha }(B_{ϵn})`$ eventually for any fixed $`r1`$ and $`ϵ>0`$, all limiting quantities are equal provided we take $`r\mathrm{}`$ and/or $`ϵ0`$ after $`n\mathrm{}`$. ∎ ## 6. Proof of main result Before we start proving our main result, Theorem 2.1, we need to ensure that if a large component is present in the graph, then it is unique. The statement we need is as follows: ###### Lemma 6.1 Let $`K_{ϵ,2}`$ be the event that $`𝒢(n,\alpha /n)`$ is either connected or has exactly two connected components, each of which is of size at least $`ϵn`$, and recall that $`K`$ is the event that $`𝒢(n,\alpha /n)`$ is connected. Then for all $`\alpha _0>0`$ and $`ϵ_0>0`$ there exists $`c_1=c_1(\alpha _0,ϵ_0)<1`$ such that for all $`ϵϵ_0`$ and all $`\alpha \alpha _0`$, $$\underset{n\mathrm{}}{lim\; sup}P_{\alpha ,n}(K^\text{c}\|K_{ϵ,2})^{1/n}<c_1.$$ (6.1) Proof. It clearly suffices to show that the ratio of $`P_{\alpha ,n}(K_{ϵ,2}K)`$ and $`P_{\alpha ,n}(K)`$ decays to zero exponentially with $`n`$, with a rate that is uniformly bounded in $`ϵϵ_0`$ and $`\alpha \alpha _0`$. In light of Theorem 2.3 and the fact that $`K_{ϵ,2}`$ only admits components that grow linearly with $`n`$, we have $$\frac{P_{\alpha ,n}(K_{ϵ,2}K)}{P_{\alpha ,n}(K)}=\text{e}^{o(n)}\underset{ϵnknϵn}{}\left(\genfrac{}{}{0pt}{}{n}{k}\right)\frac{\pi _1(\alpha k/n)^k\pi _1\left(\alpha (1k/n)\right)^{nk}}{\pi _1(\alpha )^n}\left(1\frac{\alpha }{n}\right)^{k(nk)},$$ (6.2) where $`o(n)/n`$ tends to zero uniformly in $`ϵϵ_0`$ and $`\alpha \alpha _0`$. Writing $`\varrho `$ for $`k`$$`/`$$`n`$, the expression under the sum can be bounded by $`\text{e}^{n[\mathrm{\Xi }(\varrho )\mathrm{\Xi }(0)]}`$, where $$\mathrm{\Xi }(\varrho )=S(\varrho )+\varrho \mathrm{log}\pi _1(\alpha \varrho )+(1\varrho )\mathrm{log}\pi _1\left(\alpha (1\varrho )\right)\alpha \varrho (1\varrho ).$$ (6.3) Since $`\varrho `$ is restricted to the interval $`[ϵ,1ϵ]`$, the right-hand side of (6.2) will be exponentially small if we can show $`\mathrm{\Xi }(\varrho )<\mathrm{\Xi }(0)`$ for all $`\varrho (0,1)`$ and all $`\alpha `$. As is easy to check, the function $`\varrho \mathrm{\Xi }(\varrho )`$ is symmetric about the point $`\varrho =1/2`$. Hence, if we can prove that it is strictly convex throughout $`[0,1]`$, then it is maximized at the endpoints. Introducing the function $$G(\eta )=\eta \mathrm{log}\frac{\pi _1(\eta )}{\eta }$$ (6.4) we have $$\alpha \mathrm{\Xi }(\eta /\alpha )=G(\eta )+G(\alpha \eta )+\eta (\alpha \eta ).$$ (6.5) In order to prove strict convexity of $`\mathrm{\Xi }`$, it thus suffices to show that $`G^{\prime \prime }(\eta )+1>0`$ for all $`\eta >0`$. Introducing yet another abbreviation $`q(\eta )=\eta /(1\text{e}^\eta )`$, a tedious but straightforward differentiation yields $$G^{\prime \prime }(\eta )+1=\frac{1}{q}(q^{}q)(q\text{e}^\eta 1).$$ (6.6) A direct evaluation now shows that both $`q^{}q`$ and $`q\text{e}^\eta 1`$ are negative once $`\eta >0`$. ∎ We will use the above lemma via the following simple conclusion: ###### Lemma 6.2 Let $`N_r`$ denote the number of connected components of $`𝒢(n,\alpha /n)`$ of size at least $`r`$ and let $`𝒱_r`$ be the set of vertices contained in these components. Then for all $`\alpha >0`$ and $`\varrho >ϵ>0`$ there exists $`c=c(ϵ,\varrho ,\alpha )>0`$ such that $$P_{\alpha ,n}\left(\|𝒱_{ϵn}\|=\varrho n\&N_{ϵn}=1\right)(1\text{e}^{cn})P_{\alpha ,n}\left(\|𝒱_{ϵn}\|=\varrho n\right).$$ (6.7) Proof. Clearly, (6.7) will follow if we can prove that $$P_{\alpha ,n}\left(\|𝒱_{ϵn}\|=\varrho n\&N_{ϵn}>1\right)\text{e}^{cn}P_{\alpha ,n}\left(\|𝒱_{ϵn}\|=\varrho n\right).$$ (6.8) Let $`𝒱(x)`$ denote the connected component of $`𝒢(n,\alpha /n)`$ containing the vertex $`x`$ and let $`xy`$ denote the event that $`x,y𝒱_{ϵn}`$ but $`𝒱(x)𝒱(y)=\mathrm{}`$. Then (6.8) will be proved once we show $$P_{\alpha ,n}\left(\|𝒱_{ϵn}\|=\varrho n\&xy\right)\text{e}^{2cn}P_{\alpha ,n}\left(\|𝒱_{ϵn}\|=\varrho n\right).$$ (6.9) (Indeed, the sum over $`x,y`$ adds only a multiplicative factor of $`n^2`$ on the right-hand side.) By conditioning on the set $`𝒱_{ϵn}`$ and the set $`𝒱(x)𝒱(y)`$, this inequality will in turn follow from $$P_{\alpha ,n}\left(xy\&𝒱(x)𝒱(y)=𝒱\right)\text{e}^{2cn}P_{\alpha ,n}\left(𝒱(x)𝒱(y)=𝒱\right).$$ (6.10) Indeed, let us multiply both sides by the probability that $`𝒱`$ is disconnected from the rest of the graph and that all components disjoint from $`𝒱`$ of size at least $`ϵn`$ take the total volume $`\varrho n\|𝒱\|`$. The sum over all admissible $`𝒱`$ reduces (6.10) to (6.9). We will deduce (6.10) from Lemma 6.1. Recall that $`K`$ is the event that the graph is connected and $`K_{ϵ,2}`$ is the event that it has at most two components, each of which is of size at least $`ϵn`$. We will now use these events for the restriction of $`𝒢(n,\alpha /n)`$ to $`𝒱`$: Let $`m=\|𝒱\|`$, $`\stackrel{~}{\alpha }=\alpha \frac{m}{n}`$ and $`\stackrel{~}{ϵ}=ϵ\frac{n}{m}`$. Then we have $$\{xy\}\left\{𝒱(x)𝒱(y)=𝒱\right\}K^\text{c}K_{\stackrel{~}{ϵ},2},$$ (6.11) while for the event on the right-hand side of (6.10) we simply get $$\left\{𝒱(x)𝒱(y)=𝒱\right\}=K_{\stackrel{~}{ϵ},2}.$$ (6.12) By Lemma 6.1 and the fact that $`\stackrel{~}{\alpha }\alpha `$ and $`\stackrel{~}{ϵ}ϵ`$, $$P_{\stackrel{~}{\alpha },m}(K^\text{c}\|K_{\stackrel{~}{ϵ},2})\text{e}^{c_1m},$$ (6.13) once $`n`$ is sufficiently large. But $`m2ϵn`$ and so (6.10) holds with $`c=ϵc_1`$. ∎ Now we have finally amassed all ingredients needed for the proof of our main result. Proof of Theorem 2.1. The case $`\varrho =0`$ is quickly reduced to Theorems 2.42.5 while $`\varrho =1`$ boils down to Theorem 2.3. Thus, we are down to the cases $`\varrho (0,1)`$. Let $`ϵ(0,\varrho )`$. By Lemma 6.2, we can focus on the situations with $`N_{ϵn}=1`$. To make our notation simple, let us assume that $`\varrho n`$ is an integer. Then we have $$P_{\alpha ,n}\left(\|𝒱_{ϵn}\|=\varrho n\&N_{ϵn}=1\right)=\left(\genfrac{}{}{0pt}{}{n}{\varrho n}\right)P_{\varrho n,\alpha \varrho }(K)P_{n\varrho n,\alpha (1\varrho )}(B_{ϵn})\left(1\frac{\alpha }{n}\right)^{\varrho n(1\varrho )n}.$$ (6.14) The terms on the right-hand side represent the following: the number of ways to choose the unique component of size $`\varrho n`$, the probability that this component is connected, the probability that the complement contains no component of size larger than $`ϵn`$ and, finally, the probability that the two parts of the graph do not have any edge between them. Invoking Stirling’s formula to deal with the binomial term, and plugging explicit expressions for $`P_{\varrho n,\alpha \varrho }(K)`$ and $`P_{n\varrho n,\alpha (1\varrho )}(B_{ϵn})`$ from Theorems 2.32.5, the result reduces to a simple calculation. ∎ ## Acknowledgments This research was partially supported by the NSF grants DMS-0306167, DMS-0301795 and DMS-0505356. We wish to thank anonymous referees for advice on style and literature.
warning/0506/cond-mat0506487.html	ar5iv	text	# Partly occupied Wannier functions: Construction and applications ## I Introduction A characteristic property of the single-particle eigenstates of most molecular and solid state systems is their delocalized nature. For many practical purposes this property is undesired and the construction of equivalent representations in terms of localized orbitals becomes an important issue. Within the independent-particle approximation the use of Wannier functions (WFs) allows for an exact description of the electronic groundstate in terms of a minimal set of localized orbitalswannier37 . The Wannier basis is truly minimal in the sense that the number of orbitals is just enough to accomodate the valence electrons of the system. Moreover, these localized WFs provide a formal justification of the widely used tight-binding ashcroft\_mermin and Hubbard models mahan . Being the local analogue of the extended Bloch states of solid state physics, the WFs formalize standard chemical concepts such as bonding, coordination and electron lone pairs. Among the more technical applications of Wannier functions we mention the connection to polarization theoryking-smith93 ; resta94 and their use within so-called “linear scaling” or “order-$`N`$” methods to obtain the electronic groundstategoedecker99 . Very recently numerical methods for electron transport calculations employing a Wannier function basis set have been developedcalzolari04 ; thygesen\_chemphys . In the context of molecular systems the analogue of Wannier functions for finite systems has been studied under the name ”localized molecular orbitals” boys60 ; foster60 ; edmiston63 ; pipek89 ; berghold00 ; silvestrelli98 . These are traditionally defined by an appropriate unitary transformation of the occupied single-particle eigenstates and have been used for investigation of chemical bonding. In the following we shall for simplicity use the term WF to cover also localized molecular orbitals. In 1997 Marzari and Vanderbilt developed a scheme to perform practical calculations of maximally localized Wannier functions for an isolated group of bands, i.e. a set of bands which is separated by a finite gap from all higher- and lower-lying bandsmarzari\_vanderbilt97 . Within this scheme, the usual arbitrariness inherit in the definition of the Wannier functions due to the unspecified set of unitary transformations of the Bloch states at every wave vector, is removed by requiring that the sum of second moments of the resulting WFs is minimal. The method follows the traditional idea of defining Wannier functions by a unitary transformation of the occupied (Bloch-) orbitals. In general, such methods fail to produce well localized orbitals when applied to metallic systems because the unoccupied states belonging to the partly filled valence bandsiannuzzi02 are not considered. Of course, in cases where the partly filled valence bands are separated by a gap from all higher bands, the method of Marzari and Vanderbilt still applies. However, in the more general case where the bands of interest cross and/or hybridize with other unwanted bands a different approach must be used. In this paper we demonstrate how the localization and in some cases also the symmetry of a set of WFs can be drastically improved by including selected unoccupied states in the definition of the WFs thygesen\_WFprl . The determination of the relevant unoccupied states can be viewed as a bonding-antibonding closing procedure, where occupied bonding states are paired with their antibonding counterparts to yield localized orbitals. To be more specific, consider two well-localized atomic orbitals on neighboring atoms in a molecule. If we allow the two states to hybridize, a bonding and an antibonding combination will result – combinations which may be less localized than the individual atomic orbitals. To regain the localized atomic orbitals from the molecular orbitals we need both the bonding and the antibonding combination independent of their occupation, see Fig. 1. In some cases the antibonding state may have hybridized further with other states and the state which “matches” the bonding state will be a linear combination of eigenstates. The problem we address here is the construction of a method for systematically identifying the relevant unoccupied states. We show that this can be achieved by optimizing the localization of the resulting WFs. The paper gives a more detailed and extended account of the work previously published in a Letter.thygesen\_WFprl For periodic systems the bonding-antibonding closure can be viewed as a procedure for disentangling the partly occupied valence bands from higher-lying bands. This problem has previously been addressed by Souza et al. souza01 who proposed a disentangling method based on a minimization of the change in character of the Bloch states across the Brillouin zone (BZ). While this is a natural strategy for crystalline systems, it is not clear how this disentanglement procedure applies to non-periodic systems like isolated molecules, a surface with adsorbates or a metal with impurities. The present method is related to that of Souza et al. souza01 , however, instead of minimizing the dispersion across the BZ we suggest a disentanglement procedure based exclusively on a minimization of the spread of the WFs. In this way we omit any reference to the wave-vector and are therefore not limited to periodic systems. The generality of the method is demonstrated by application to three different systems: an isolated $`\text{Si}_5`$ cluster, a copper crystal, and a Cu(100) surface with nitrogen adsorbed. Our results for the copper crystal are very similar to those obtained by Souza et al. souza01 , and this indicates the similarity of the two localization schemes for periodic systems. The paper is organized as follows: In Sec. II we introduce the spread functional and outline the strategy behind the localization algorithm. In Sec. III we give the formal definition of partly occupied WFs in the limiting case of a large supercell and derive the corresponding expressions for the gradient of the spread functional. The extension to periodic systems is discussed in Sec. IV. In Sec. V we apply the method to a $`\text{Si}_5`$ cluster, a copper crystal and a Cu(100) surface with adsorbed nitrogen. ## II Description of the method In this section we introduce the spread functional used to measure the degree of localization of a set of orbitals, and give an introductory description of the localization scheme including its relation to the method of Souza et al. souza01 ### II.1 Spread functional Within the localization scheme of Marzari and Vanderbilt marzari\_vanderbilt97 the spread of a set of functions $`\{w_n(𝕣)\}_{n=1}^N`$ is measured by the sum of second moments $$S=\underset{n=1}{\overset{N}{}}(w_n\|r^2\|w_nw_n\|𝕣\|w_n^2).$$ (1) When periodic boundary conditions are applied, as in the present study, and the supercell is sufficiently large, the minimization of $`S`$ is equivalent to the maximization of resta99 $$\mathrm{\Omega }=\underset{n=1}{\overset{N}{}}\underset{\alpha =1}{\overset{N_G}{}}W_\alpha \|Z_{\alpha ,nn}\|^2,$$ (2) where the matrix $`Z_\alpha `$ is defined as $$Z_{\alpha ,nm}=w_n\|e^{i𝔾_\alpha 𝕣}\|w_m.$$ (3) The $`\{𝔾_\alpha \}`$ is a set of at most six reciprocal lattice vectors and $`\{W_\alpha \}`$ are corresponding weights which account for the shape of the unit cell. For a definition and discussion of these quantities we refer to Refs. silvestrelli98, ; berghold00, . ### II.2 The localization scheme The starting point is the set of single-particle eigenstates, $`\{\psi _n\}`$, resulting from a conventional electronic structure calculation. For simplicity we shall assume that the system is isolated or is contained in a large supercell such that reference to $`𝕜`$-points can be omitted. The aim is to obtain a set of $`N_w`$ localized WFs with the property that any eigenstate below a specified energy, $`E_0`$, can be exactly reproduced as a linear combination of the WFs. An obvious way to achieve this would be to apply the method of Marzari and Vanderbilt to compute the unitary transformation of the $`N_w`$ lowest eigenstates leading to the most localized WFs. The problem with this strategy is, however, that it is in general not possible to localize all WFs simultaneously, and the problem cannot be overcome by increasing $`N_w`$. Instead, we define an external localization space as the space spanned by the $`N_b`$ lowest-lying eigenstates ($`N_b>N_w`$). Within this space we consider the subspace spanned by the eigenstates with energy below $`E_0`$, together with $`L`$ extra degrees of freedom (EDF). We shall refer to this subspace as the active localization space or simply the localization space. The EDF are assumed to be orthogonal and $`L`$ is chosen such that the dimension of the active localization space equals $`N_w`$. We then perform a simultaneous optimization of the WFs within the active localization space and of the active localization space itself. In practice this is achieved by optimizing an $`N_w\times N_w`$ unitary matrix together with the coordinates of the EDF such that the functional $`\mathrm{\Omega }`$ becomes maximal. It is the determination of the EDF that distinguish our method from that of Souza et al. souza01 In the latter, the spread functional is decomposed into two terms: $`\mathrm{\Omega }=\mathrm{\Omega }_I+\stackrel{~}{\mathrm{\Omega }}`$, where $`\mathrm{\Omega }_I`$ is related to the $`k`$-space dispersion of the band-projection operator, see Ref. souza01, . In the first step, the EDF are determined by maximizing $`\mathrm{\Omega }_I`$, which depends only on the localization space itself and not on the internal unitary transformation. In the second step $`\stackrel{~}{\mathrm{\Omega }}`$, or equivalently $`\mathrm{\Omega }`$, is then maximized within the fixed localization space. It is clear, that the separate maximization of $`\mathrm{\Omega }_I`$ and $`\stackrel{~}{\mathrm{\Omega }}`$ does not amount to the global maximization of $`\mathrm{\Omega }`$ that we propose here. We shall, however, see that the two methods lead to very similar results in the case of periodic systems. ## III Large supercells In this section we give a detailed description of the localization scheme in the limiting case of a large supercell where a $`\mathrm{\Gamma }`$-point sampling of the first Brillouin zone is a good approximation. For simplicity we discuss this case separately before extending it to periodic systems, although the latter contains the former as a special case. After giving the definition of partly occupied Wannier functions we derive expressions for the gradients of the spread functional and discuss how to combine these with a Lagrange multiplier scheme to determine the maximum of $`\mathrm{\Omega }`$. ### III.1 Definition of partly occupied Wannier functions We denote the total number of eigenstates obtained from the electronic structure calculation by $`N_b`$ and the number of eigenstates below the energy $`E_0`$ by $`M`$. Our aim is to construct a set of $`N_w`$ WFs which span at least the $`M`$ lowest-lying eigenstates. The remaining $`L=N_wM`$ degrees of freedom are simply used to improve the localization of the resulting WFs as much as possible. We expand the WFs in terms of the $`M`$ lowest lying eigenstates and $`L`$ extra degrees of freedom, $`\{\varphi _l\}`$, belonging to the $`(N_bM)`$-dimensional space of eigenstates with energy above $`E_0`$: $$w_n=\underset{m=1}{\overset{M}{}}U_{mn}\psi _m+\underset{l=1}{\overset{L}{}}U_{M+l,n}\varphi _l,$$ (4) where the extra degrees of freedom (EDF) are written as $$\varphi _l=\underset{m=1}{\overset{N_bM}{}}c_{ml}\psi _{M+m}.$$ (5) The columns of the matrix $`c`$ are orthonormal and represent the coordinates of the EDF with respect to the eigenstates lying above $`E_0`$. The matrix $`U`$ is unitary and represents a rotation of the functions $`\{\psi _1,\mathrm{},\psi _M,\varphi _1,\mathrm{},\varphi _L\}`$. In order to simplify the notation we introduce the matrices $$C=\left[\begin{array}{cc}I^{M\times M}& 0\\ 0& c\end{array}\right],V=CU=\left[\begin{array}{c}U^M\\ cU_L\end{array}\right],$$ (6) where $`U^M`$ and $`U_L`$ denotes the $`M`$ uppermost and $`L`$ lowermost rows of $`U`$, respectively. The $`i`$th column of $`V`$ gives the coordinates of $`w_i`$ with respect to the full set of eigenstates $`\{\psi _n\}`$. Substituting the expansions (4) and (5) into Eq. (3) we obtain a compact matrix expression $$Z_\alpha =V^{}Z_\alpha ^{(0)}V=U^{}C^{}Z_\alpha ^{(0)}CU,$$ (7) where $`Z_\alpha ^{(0)}`$ is obtained from Eq. (3) by using the eigenstates $`\{\psi _n\}`$ in the inner product, $$Z_{\alpha ,nm}^{(0)}=\psi _n\|e^{i𝔾_\alpha 𝕣}\|\psi _m.$$ (8) ### III.2 Gradient of $`\mathrm{\Omega }`$ Through Eq. (7) the spread functional, $`\mathrm{\Omega }`$, in Eq. (2) becomes a function of the matrices $`U`$ and $`c`$. The maximum of $`\mathrm{\Omega }`$ can be found iteratively by updating $`U`$ and $`c`$ in the direction given by the gradient. In the following we derive expressions for the gradient of $`\mathrm{\Omega }`$. We write the unitary matrix at iteration $`n`$ as $`U^{(n)}=U^{(n1)}\mathrm{exp}(A)`$, where $`A`$ is an anti-hermitian matrix. Since we are only concerned with small variations, we expand the exponential to first order, i.e. $`\mathrm{exp}(A)1A`$. Inserting this into Eqs. (7) and (2) we find $$\frac{\mathrm{\Omega }}{A_{ij}}=\underset{\alpha =1}{\overset{N_G}{}}W_\alpha [Z_{\alpha ,ji}(Z_{\alpha ,jj}^{}Z_{\alpha ,ii}^{})Z_{\alpha ,ij}^{}(Z_{\alpha ,ii}Z_{\alpha ,jj})].$$ (9) All matrices in this expression refer to iteration $`n1`$. The new rotation at iteration $`n`$ is then obtained by multiplying $`U^{(n1)}`$ by $`\mathrm{exp}[d(_A\mathrm{\Omega })]`$ where $`d`$ is the length of the steepest-ascent step and $`[_A\mathrm{\Omega }]_{ij}=\mathrm{\Omega }/A_{ij}`$. We now turn to the problem of determining the steepest uphill direction of $`\mathrm{\Omega }`$ with respect to variations in $`c`$. In general, for a real-valued function $`f(z=x+iy)`$ the direction of steepest ascent with respect to $`z`$ is given by $$\frac{f}{z^{}}\frac{1}{2}(\frac{f}{x}+i\frac{f}{y}).$$ (10) To calculate the gradient $`\mathrm{\Omega }/c_{ij}^{}`$ we use that $$\frac{\|Z_{\alpha ,nn}\|^2}{c_{ij}^{}}=Z_{\alpha ,nn}\frac{Z_{\alpha ,nn}^{}}{c_{ij}^{}}+Z_{\alpha ,nn}^{}\frac{Z_{\alpha ,nn}}{c_{ij}^{}}.$$ (11) From Eq. (7) it follows that $`{\displaystyle \frac{Z_{\alpha ,nn}}{c_{ij}^{}}}`$ $`=`$ $`{\displaystyle \underset{abcd}{}}U_{na}^{}{\displaystyle \frac{C_{ab}^{}}{c_{ij}^{}}}Z_{\alpha ,bc}^{(0)}C_{cd}U_{dn}`$ $`+`$ $`{\displaystyle \underset{abcd}{}}U_{na}^{}C_{ab}^{}Z_{\alpha ,bc}^{(0)}{\displaystyle \frac{C_{cd}}{c_{ij}^{}}}U_{dn},`$ and from definition (6) $`{\displaystyle \frac{C_{nm}}{c_{ij}^{}}}`$ $`=`$ $`0`$ (13) $`{\displaystyle \frac{C_{nm}^{}}{c_{ij}^{}}}`$ $`=`$ $`\delta _{m,M+i}\delta _{n,M+j}.`$ (14) It is now easy to establish that $`{\displaystyle \frac{Z_{\alpha ,nn}}{c_{ij}^{}}}=[Z_\alpha ^{(0)}V]_{M+i,n}U_{M+j,n}^{}`$ (15) $`{\displaystyle \frac{Z_{\alpha ,nn}^{}}{c_{ij}^{}}}=[(Z_\alpha ^{(0)})^{}V]_{M+i,n}U_{M+j,n}^{}.`$ (16) Combining Eq. (11) with (15) and (16) we arrive at the desired expression $$\frac{\mathrm{\Omega }}{c_{ij}^{}}=\underset{\alpha =1}{\overset{N_G}{}}W_\alpha [Z_\alpha ^{(0)}V\text{D}(Z_\alpha ^{})U^{}+(Z_\alpha ^{(0)})^{}V\text{D}(Z_\alpha )U^{}]_{M+i,M+j},$$ (17) where $`\text{D}(Z_\alpha )`$ is a diagonal matrix with $`(Z_{\alpha ,nn})`$ in the diagonal. To treat the constraint that the EDF $`\{\varphi _l\}`$ should be orthonormal during the maximization procedure we introduce the Lagrange multipliers $`\lambda _{ij}`$ and perform an unconstrained maximization of the functional $$\mathrm{\Omega }_L=\mathrm{\Omega }\underset{ij}{}\lambda _{ij}\varphi _i\|\varphi _j.$$ (18) The Lagrange multipliers are initially unknown and must be estimated at each iteration. At the maximum we have $`_c^{}\mathrm{\Omega }_L=0`$ which is equivalent to the condition $$_c^{}\mathrm{\Omega }c\text{ }\lambda ^\text{T}=0.$$ (19) Multiplying by $`c^{}`$ from the left leads to $$\lambda ^\text{T}=c^{}_c^{}\mathrm{\Omega }.$$ (20) This relation can be used to estimate the Lagrange multipliers at each iteration. A step of length $`d`$ in the steepest uphill direction is thus accomplished by adding to $`c`$ the matrix $`d(1cc^{})_c^{}\mathrm{\Omega }`$, followed by an orthonormalization of the columns of $`c`$. ## IV Periodic systems We consider a periodic system with a unit cell defined by basis vectors $`𝕒_1,𝕒_2,𝕒_3`$ which in turn define the basis vectors of the reciprocal lattice $`𝕓_1,𝕓_2,𝕓_3`$. The Bloch states, $`\{\psi _{n𝕜}\}`$, resulting from the electronic structure calculation are characterized by a band index $`n`$ and a crystal momentum $`𝕜`$. The total number of bands is denoted by $`N_b`$ and the number of eigenstates at a given $`𝕜`$-point with energy below $`E_0`$ is denoted by $`M_𝕜`$. We assume a uniform sampling of the first BZ such that any $`𝕜`$-point can be written as $$𝕜=\frac{n_1}{N_1}𝕓_1+\frac{n_2}{N_2}𝕓_2+\frac{n_3}{N_3}𝕓_3,$$ (21) where $`N_i`$ is the number of $`𝕜`$-points in the direction $`𝕓_i`$ and $`n_i=0,\mathrm{},N_i1`$. Note that the $`\mathrm{\Gamma }`$ point is always included. With this convention the Bloch states, $`\{\psi _{n𝕜}\}`$, correspond exactly to the $`\mathrm{\Gamma }`$-point eigenstates of the repeated cell defined by the extended basis vectors $`N_1𝕒_1,N_2𝕒_2,N_3𝕒_3`$. An alternative way of stating this correspondence is to say that the $`𝕜`$-points in Eq. (21) fall on the reciprocal lattice of the repeated cell, see Fig. 2. As we shall see below, this correspondence allows us to use the spread functional $`\mathrm{\Omega }`$ defined in Eq. (2) also for the periodic system. We stress that the formalism developed in the following section contains the $`\mathrm{\Gamma }`$-point formalism described in the preceding sections as a special case. ### IV.1 Definition of partly occupied Wannier functions We write the $`n`$th Wannier function related to unit cell $`i`$ as $$w_{i,n}=\frac{1}{\sqrt{N_k}}\underset{𝕜}{}e^{i𝕜_i}\stackrel{~}{\psi }_{n𝕜},$$ (22) where $`N_k`$ is the total number of $`𝕜`$-points and $`\stackrel{~}{\psi }_{n𝕜}`$ is a generalized Bloch state to be defined below marzari\_vanderbilt97 . Each generalized band, i.e. each set $`\{\stackrel{~}{\psi }_{n𝕜}\}`$ for fixed $`n`$, gives rise to one WF per unit cell. These WFs are simply related by translation, i.e. $`w_{i,n}(𝕣)=w_{0,n}(𝕣_i)`$, and thus it suffices to consider the WFs of the cell at the origin. In doing this we can omit the cell index and simply denote the WFs by $`\{w_n\}`$. We denote the number of WFs per cell by $`N_w`$. Following the idea behind Eq. (4) we expand the generalized Bloch state $`\stackrel{~}{\psi }_{n𝕜}`$ in terms of the $`M_𝕜`$ lowest lying Bloch states and $`L_𝕜`$ extra degrees of freedom, $`\{\varphi _{l𝕜}\}`$, from the remaining $`(N_bM_𝕜)`$-dimensional space $$\stackrel{~}{\psi }_{n𝕜}=\underset{m=1}{\overset{M_𝕜}{}}U_{mn}^𝕜\psi _{m𝕜}+\underset{l=1}{\overset{L_𝕜}{}}U_{M_𝕜+l,n}^𝕜\varphi _{l𝕜},$$ (23) where the EDF are expanded as $$\varphi _{l𝕜}=\underset{m=1}{\overset{N_bM_𝕜}{}}c_{ml}^𝕜\psi _{M_𝕜+m,𝕜}.$$ (24) The number of EDF at a given $`𝕜`$-point is determined by the condition $`L_𝕜+M_𝕜=N_w`$. If $`M_𝕜`$ exceeds $`N_w`$, we simply put $`M_𝕜=N_w`$. Due to the exact correspondence between the Bloch states $`\{\psi _{n𝕜}\}`$ and the $`\mathrm{\Gamma }`$-point eigenstates of the repeated cell, we can use the functional (2) to measure the spread of the Wannier functions. The matrices $`Z_\alpha `$ are still defined by Eq. (3) but it should be remembered that the inner product as well as the reciprocal lattice vector $`𝔾_\alpha `$ now refer to the repeated cell. From Eqs. (23,24) we find the following generalization of Eq. (7) $$Z_\alpha =\underset{𝕜,𝕜^{}}{}Z_\alpha ^{𝕜𝕜^{}},$$ (25) where $$Z_\alpha ^{𝕜𝕜^{}}=(U^𝕜)^{}(C^𝕜)^{}Z_\alpha ^{(0),𝕜𝕜^{}}C^𝕜^{}U^𝕜^{}.$$ (26) The matrix $`C^𝕜`$ is given by the obvious $`𝕜`$-point analogue of Eq. (6) and the matrix $`Z_\alpha ^{(0),𝕜𝕜^{}}`$ is defined by $$Z_{\alpha ,nm}^{(0),𝕜𝕜^{}}=\psi _{n𝕜}\|e^{i𝔾_\alpha 𝕣}\|\psi _{m𝕜^{}}.$$ (27) Most of the matrices $`Z_\alpha ^{(0),𝕜𝕜^{}}`$ are in fact zero. Writing the Bloch functions as $`\psi _{n𝕜}=u_{n𝕜}(𝕣)\mathrm{exp}(i𝕜𝕣)`$, where $`u_{n𝕜}`$ has the periodicity of the lattice, we get $$Z_{\alpha ,nm}^{(0),𝕜𝕜^{}}=u_{n𝕜}^{}(𝕣)u_{m𝕜^{}}(𝕣)e^{i(𝕜^{}𝕜𝔾_\alpha )𝕣}\text{d}𝕣,$$ (28) which is non-zero only when $$𝕜^{}=𝕜+𝔾_\alpha .$$ (29) Here it is implicit that $`𝕜`$ and $`𝕜^{}`$ belong to the first BZ and thus it might be necessary to translate $`𝕜^{}`$ by a reciprocal lattice vector. The relation between $`𝕜`$ and $`𝕜^{}`$ is illustrated in Fig. 2. Note that the condition in Eq. (29) reduces the double sum in Eq. (25) to a single sum over $`𝕜`$. The derivation of the gradient of $`\mathrm{\Omega }`$ follows closely the $`\mathrm{\Gamma }`$-point case discussed in Sec. III.2 and is therefore omitted. The result is $$\frac{\mathrm{\Omega }}{A_{ij}^𝕜}=\underset{\alpha =1}{\overset{N_G}{}}W_\alpha [(Z_{\alpha ,jj})^{}Z_{\alpha ,ji}^{𝕜𝔾_\alpha ,𝕜}+Z_{\alpha ,jj}(Z_{\alpha ,ij}^{𝕜,𝕜+𝔾_\alpha })^{}(Z_{\alpha ,ii})^{}Z_{\alpha ,ji}^{𝕜,𝕜+𝔾_\alpha }Z_{\alpha ,ii}(Z_{\alpha ,ij}^{𝕜𝔾_\alpha ,𝕜})^{}].$$ (30) $$\frac{\mathrm{\Omega }}{(c_{ij}^𝕜)^{}}=\underset{\alpha =1}{\overset{N_G}{}}W_\alpha [Z_\alpha ^{(0),𝕜,𝕜+𝔾_\alpha }V^{𝕜+𝔾_\alpha }\text{D}(Z_\alpha ^{})(U^𝕜)^{}+(Z_\alpha ^{(0),𝕜𝔾_\alpha ,𝕜})^{}V^{𝕜𝔾_\alpha }\text{D}(Z_\alpha )(U^𝕜)^{}]_{M_𝕜+i,M_𝕜+j}.$$ (31) We note that these expressions, of course, reduce to Eqs. (9,17) in the limit of a single $`𝕜`$-point. The maximization of $`\mathrm{\Omega }`$ proceeds along the same lines as for the $`\mathrm{\Gamma }`$-point case, except that Lagrange multipliers are needed for each $`𝕜`$-point. For example the analogue of Eq. (18) reads $$\mathrm{\Omega }_L=\mathrm{\Omega }\underset{ij,𝕜}{}\lambda _{ij,𝕜}\varphi _{i𝕜}\|\varphi _{j𝕜}.$$ (32) ### IV.2 Optimizing the number of extra degrees of freedom For given values of $`N_b`$, $`N_w`$ and $`E_0`$, the algorithm introduced above produces the $`N_w`$ most localized WFs that can be formed within the external localization space when all eigenstates below $`E_0`$ should be exactly reproducible in terms of the WFs. It remains to determine the optimal values for $`N_b`$ and $`N_w`$ for a given $`E_0`$. Let us start by considering the situation where $`N_b`$ has been fixed at a value which is large enough to include all anti-bonding states relevant for the localization. In practice this typically means $`10`$ eV above the Fermi level. It seems as a natural strategy to choose $`N_w`$ such that the localization per orbital is maximal. To quantify this condition we define the average localization per orbital as $$\mathrm{\Omega }=\frac{\mathrm{\Omega }[E_0,N_b,N_w]}{N_w},$$ (33) where we have indicated the dependence of $`\mathrm{\Omega }`$ on the three parameters explicitly. We note that since the value of $`\mathrm{\Omega }`$ also depends on the size and shape of the supercell, it does not make sense to compare the value of $`\mathrm{\Omega }`$ for systems described in different supercells. Fixing $`N_w`$ on the basis of $`\mathrm{\Omega }`$ represents a completely general criterion which can be applied in any situation. However, the localization procedure must be carried out for several values of $`N_w`$ which might be a tedious task depending on the size of the system. We next consider the situation when $`N_b`$ is also allowed to change. Formally, the global maximum of $`\mathrm{\Omega }`$ is attained in the limit where both $`N_b`$ and $`N_w`$ tend to infinity in which case an infinite set of completely localized delta functions can be realized. However, we have found that for practical values of $`N_b`$ where very high energy states are not included in the external localization space, $`\mathrm{\Omega }`$ will have a local maximum for some $`N_w`$, and the position of the maximum is not sensitive to the actual value of $`N_b`$. Thus, it is indeed possible to determine an optimal value of $`N_w`$ by maximizing $`\mathrm{\Omega }`$. Alternatively it is often possible to determine a value for $`N_w`$ based on symmetry arguments, chemical intuition, or a closed band condition. As we shall see in the following examples the two criteria for determining $`N_w`$ lead to similar results. ### IV.3 Start guess for $`U^𝕜`$ and $`c^𝕜`$ For small systems we have found that the localization algorithm is quite stable and usually leads to the global maximum independently of the initial value of the matrices $`\{U^𝕜\},\{c^𝕜\}`$. For larger systems, however, there is a risk of getting stuck in a local maximum and in such cases the start guess becomes important. It is then natural to start from a set of simple orbitals located either at the atoms or at the bond centers. Let $`\{f_\nu \}`$ denote such a set of simple orbitals. The question is how to transform this into the matrices $`\{U^𝕜\},\{c^𝕜\}`$. To this end we project the initial orbitals onto the subspace spanned by the Bloch states at each $`𝕜`$-point: $$\stackrel{~}{f}_{\nu 𝕜}=\underset{n=1}{\overset{N_b}{}}\psi _{n𝕜}\|f_\nu \psi _{n𝕜}.$$ (34) The following procedure is carried out for each $`𝕜`$-point separately. For fixed $`𝕜`$ we regard $`\psi _{n𝕜}\|f_\nu `$ as a matrix in the indices $`n,\nu `$. Its columns represent the coordinates of the $`\stackrel{~}{f}_{𝕜\nu }`$ with respect to the Bloch states $`\{\psi _{n𝕜}\}_{n=1}^{N_b}`$ and as such it is a (non-orthogonalized) version of the matrix $`V^𝕜`$, see Eq. (6). After a normalization of the columns of $`\psi _{n𝕜}\|f_\nu `$ we compute the norm of the component of $`\stackrel{~}{f}_{\nu 𝕜}`$ orthogonal to the occupied subspace: $$\stackrel{~}{f}_{\nu 𝕜}^{}^2=\underset{n=M(𝕜)}{\overset{N_b}{}}\|\psi _{n𝕜}\|f_\nu \|^2.$$ (35) The first EDF is chosen as a normalized version of the $`\stackrel{~}{f}_{𝕜\nu }^{}`$ for which $`\stackrel{~}{f}_{𝕜\nu }^{}`$ is the largest. The remaining $`\stackrel{~}{f}_\nu ^{}`$’s are then orthogonalized onto this vector and the process is repeated until all EDF, and thus $`c^𝕜`$, have been determined. Finally the identity $`U^𝕜=(C^𝕜)^{}V^𝕜`$ with $`V^𝕜\psi _{n𝕜}\|f_\nu `$ determines $`U^𝕜`$. Since the $`\stackrel{~}{f}_{\nu 𝕜}`$ are not necessarily orthogonal, the columns of the resulting $`U^𝕜`$ must be explicitly orthogonalized. ## V Results In the following sections we apply the localization scheme to three different systems. To demonstrate the generality of the method we consider both isolated and metallic systems as well as a metal surface with adsorbed impurities. In Sec. V.1 we construct partly occupied WFs for an isolated $`\text{Si}_5`$ cluster and illustrate how different sets of WFs can be obtained by varying the number of extra degrees of freedom. In Sec. V.2 we investigate the WFs of a Cu(fcc) crystal and compare the results with those obtained by Souza and co-workers souza01 who studied the same system using a different but related method. Finally, in Sec. V.3 we perform a detailed WF analysis for a Cu(100) surface with 0.5 mono-layers of nitrogen. In all calculations we use a plane-wave based DFT code dacapo to obtain the Kohn-Sham eigenstates, and we describe the ion potential by Vanderbilt ultrasoft pseudopotentials vanderbilt90 . To ensure a proper convergence of the unoccupied states employed in the localization scheme, the DFT calculations have been converged with respect to the full set of Kohn-Sham eigenvalues. In the appendix we explain how to extend the localization scheme to ultrasoft pseudopotentials. ### V.1 $`\text{Si}_5`$ cluster As an example of an isolated system we consider an $`\text{Si}_5`$ cluster in its ground-state geometry raghavachari85 , see Fig. 3(a). We use a cubic supercell of length 16 Å and sample the first BZ at the $`\mathrm{\Gamma }`$-point. To test the dependence on the size of the external localization space we consider the two cases $`N_b=30`$ and $`N_b=100`$. We set $`M=10`$ corresponding to the number of occupied states, and calculate the average localization per WF, $`\mathrm{\Omega }=\mathrm{\Omega }/N_w`$, for $`L=0,\mathrm{},7`$. The result is shown in Fig. 4. For $`L=0`$ there is no difference between the two cases since the WFs are constructed entirely from the occupied eigenstates. However, for $`L0`$ the larger space available for the extra degrees of freedom leads to an improved localization when $`N_b=100`$. Apart from this general improvement in localization, there is no qualitative difference between the WFs obtained with $`N_b=30`$ and $`N_b=100`$ for a given $`L`$. We note that both curves have a maximum for $`L=4`$, corresponding to a total of 14 WFs. This particular set of WFs together with their centers is shown in Fig. 3(b). The fact that this set of WFs respect the symmetry of the cluster is a special property of the $`L=4`$ solution: for other values of $`L`$, including $`L=0`$, the WFs break the symmetry of the $`\text{Si}_5`$ cluster. This indicates that the solution corresponding to the maximal value of $`\mathrm{\Omega }`$ has a special meaning. Indeed, the value $`N_w=14`$ could also have been anticipated from physical arguments. Starting from a set of four $`sp^3`$ orbitals located at each Si atom we expect bonding and anti-bonding states to form between pairs of aligned orbitals belonging to nearest neighbor pairs of Si atoms. These bonding states can be identified as the six bond-centered WFs shown to the far left in Fig. 3. The two “top” Si atoms have three nearest neighbors and thus a single $`sp^3`$ orbital is left as a lone pair (middle WF). The remaining three Si atoms each have two nearest neighbors and consequently two $`sp^3`$ orbitals are left as lone pairs (rightmost WF). In total this adds up to 14 orbitals. The anti-bonding counterparts of the bonding states formed between nearest neighbors are not brought into play for $`L=4`$, because they are much less localized than the bonding states. However, by setting $`L=10`$ and thus searching for a total of 20 WFs, the anti-bonding states are picked out as EDF and we obtain a full set of $`sp^3`$ orbitals. This solution has, however, a smaller value for $`\mathrm{\Omega }`$ than the solution at $`L=4`$. ### V.2 Copper crystal To illustrate the method in the case of a periodic system we consider the construction of WFs for a copper crystal. This system was also studied by Souza et al.souza01 using their disentangling method to obtain the WFs. Our results are in close agreement with those obtained by Souza et al., and this indicates the similarity of the two methods for periodic systems. We use the primitive fcc unit cell and sample the first BZ on a uniform (11,11,11) Monckhorst pack grid containing the $`\mathrm{\Gamma }`$-point. To obtain a minimal set of WFs describing the Cu $`d`$\- and $`s`$-bands we set $`N_w=6`$. We construct two sets of WFs corresponding to two different values of $`E_0`$: (i) $`E_0=0.0`$ eV and (ii) $`E_0=3.0`$ eV, relative to the Fermi level. In the first case the resulting WFs will span at least the occupied subspace and thus the electronic structure described by the WFs will be correct below $`E_F`$. In the second case the electronic structure will be correct up to 3 eV above $`E_F`$, however, since this is a stronger restriction on the localization space we must expect that the resulting WFs will be less localized than those obtained in (i). In Figs. 5 and 6 we show the original DFT bands together with the approximate bands computed by diagonalizing the Hamiltonian within the subspace spanned by the WFs of case (i) and (ii), respectively. In both cases we see a very good agreement between the exact and approximate bands below $`E_0`$. At higher energies the approximate bands deviate from the exact bands, indicating that the EDF which optimize the localization of the WFs do not coincide with specific Bloch eigenstates. The quality of the WF bands below $`E_0`$ depends on the number of $`𝕜`$-points used to construct the WFs. This is because the band diagram must be constructed starting from fully localized functions, which means that the coupling matrix elements must be truncated beyond a cut-off distance given approximately by $`N_i/2`$ unit cells in the direction $`𝕒_i`$. Thus the repeated cell, or equivalently the number of $`𝕜`$-points, must be so large that the WFs have decayed sufficiently between the repeated images. Both sets of WFs consist of five atom-centered $`d`$-orbitals and a single $`s`$-like orbital centered in one of the two interstitial sites. The $`d`$-orbitals are more or less identical for the two cases, and two examples are shown in Fig. 7. Contour plots of the $`s`$-like orbital is shown in Fig. 8 (b) and (c) for case (i) and (ii), respectively. The plots indicate that the $`s`$-orbital of case (ii) is less localized than the one obtained in case (i). That this is indeed correct follows from the value of the spread functional, $`\mathrm{\Omega }`$, which is higher for (i) than for (ii). The minimal set of WFs obtained with $`N_w=6`$ breaks the symmetry of the fcc crystal because the $`s`$-like orbital is located in one of the interstitial sites leaving the other empty. As demonstrated by Souza et al. the symmetry can be restored by using seven WFs per primitive cell instead of six. In Fig. 9 we show the band structure obtained from a set of WFs generated with $`N_w=7`$ and $`E_0=0.0`$ eV. We note that very high-energetic states are now selected as the optimal EDF. This solution can therefore only be obtained for rather large external localization spaces, i.e. $`N_b9`$. The five $`d`$-like WFs are unchanged, but now we obtain two equivalent $`s`$-like WFs located in each of the two interstitial sites thereby restoring the fcc symmetry, see Fig. 8(d). We have calculated the average localization $`\mathrm{\Omega }`$ for $`N_w=6,7,8`$, and found that the maximum is attained for the symmetric solution with $`N_w=7`$. ### V.3 Nitrogen absorption on Cu(100) In this section we study the WFs of a copper (100)-surface covered with half a mono-layer of nitrogen atoms. As the system is neither periodic (in all directions) nor isolated, it represents a very general situation. The section is divided into two parts. In the first part the WFs are constructed and analyzed, and in the second part we use the obtained WFs to study the chemisorption of nitrogen within the Newns Anderson model. #### V.3.1 Wannier function analysis We model the Cu(100) surface by a slab with a thickness of two atomic layers. The supercell contains four Cu atoms and a single N atom adsorbed in a hollow site, and its height is such that the surface slabs are separated by 9.0 Å of vacuum. A topographic top-view of the surface is shown in Fig. 10. We sample the first BZ on a uniform (7,7,1) Monckhorst pack grid containing the $`\mathrm{\Gamma }`$-point. Let us start by considering what we can expect to find on the basis of our previous experience. First, the result from the copper crystal suggests that a minimal description of the metal surface is obtained with five $`d`$-orbitals and an $`s`$-like orbital per Cu atom. Since there are four Cu atoms per supercell this gives a total of 24 WFs. Next, the similarity between the valency of N and Si together with our experience from the $`\text{Si}_5`$ cluster points to a description of the nitrogen atom in terms of $`sp^x`$-hybrides. In Fig. 11 we have plotted the average localization, $`\mathrm{\Omega }`$, of the obtained WFs as a function of $`N_w`$ for three different sizes of the external localization space corresponding to $`N_b=35,40,50`$. In all cases we have set $`E_0=E_F`$ in order to ensure that the occupied eigenstates are exactly reproduced by the WFs. As expected, the localization improves as the size of the external localization space increases. In addition, the maximum of $`\mathrm{\Omega }`$ shifts towards larger $`L`$-values as $`N_b`$ is increased. Specifically the maximum shifts from $`N_w=27`$ to $`N_w=29`$ as $`N_b`$ is increased from 35 to 50. This is not unexpected since we know that $`\mathrm{\Omega }`$ will be a monotonically increasing function of $`L`$ in the limit $`N_b\mathrm{}`$, see discussion in Sec. IV.2. Again we stress that it is only the degree of localization of the WFs that change with $`N_b`$ for a fixed $`L`$, and not their qualitative form. Thus the chemical picture provided by the WFs does not change with $`N_b`$. In fact, for all the values of $`N_b`$ we obtain 20 highly localized $`d`$-orbitals (five located on each of the four Cu atoms) and four $`sp^3`$ orbitals centered on the N atom, see Fig. 13. The remaining $`N_w24`$ WFs are the less localized $`s`$-like orbitals of Cu. Thus, as $`N_w`$ is increased beyond 24, the number of $`s`$-like Cu WFs simply increases correspondingly. To gain further insight into the dependence of the WFs on $`N_b`$ and $`N_w`$, we show in Fig. 12 the average localization of the $`d`$-, $`sp^3`$\- and $`s`$-orbitals, separately. It is clear that the $`N_b`$-dependence as well as the maximum of $`\mathrm{\Omega }`$ are almost exclusively related to the Cu $`s`$-orbitals. Except for the case $`N_b=50`$, which is in fact somewhat extreme since states of 20 eV above the Fermi level are included in the external localization space, the average spread of the Cu $`s`$-orbitals is maximal for $`N_w=28`$. This corresponds to one $`s`$-orbital per Cu which is exactly what we anticipated from the analysis of the copper crystal. We end by summarizing the chemical picture obtained from the WF analysis: For $`N_w=28`$ the Cu surface is described by the minimal set of WFs consisting of five $`d`$\- and one $`s`$-like orbital per atom. For the nitrogen we obtain four $`sp^3`$ hybrids oriented as indicated in Fig. 13. #### V.3.2 Adsorption in the Newns Anderson model The WFs can be used to obtain a detailed and consistent picture of the hybridization occurring between the nitrogen states and the states of the substrate. As we shall see the analysis gives a complete account for the shape of the projected density of states of a given N orbital, in terms of the bare orbital energy, a coupling strength, and the density of states of the so-called group orbital. In the Newns Anderson model, one considers an adsorbate state, $`\|a`$, of energy $`\epsilon _a=a\|H\|a`$, coupled to a continuum of states, $`\|k`$, representing the substrate. The coupling matrix elements are denoted by $`V_k=a\|H\|k`$. A particularly useful formulation can be obtained by introducing the normalized group orbital, $`\|g=V^1_kV_k\|k`$, where $`V=(_k\|V_k\|^2)^{1/2}`$. It is easily checked, that the coupling between $`\|a`$ and any substrate state orthogonal to $`\|g`$ vanishes. Consequently $`\|a`$ is coupled to the substrate via the group orbital only, and the coupling is given by $`V`$, i.e. $`V=a\|H\|g`$. Physical quantities such as the projected density of states (PDOS) of the adsorbate state and the hybridization part of the adsorption energy, can be obtained from the retarded adsorbate Green’s function, which in turn follows from the three quantities $`\epsilon _a`$, $`V`$ and $`\rho _g^0(\epsilon )`$, where $`\rho _g^0`$ denotes the PDOS of the group orbital in the absence of coupling to the adsorbate state. Often $`\rho _g^0(\epsilon )`$ is referred to as the band to which the adsorbate is coupled. The $`sp^3`$ WFs of the N atom are not well suited as a starting point for applications of the Newns Anderson model, since they do not represent the energy levels of the free atom. This problem can be overcome by diagonalizing the Hamiltonian matrix in the WF basis, within the subspace spanned by the four $`sp^3`$ orbitals. The result of the subspace diagonalization is a set of four atomic orbitals consisting of one $`s`$-like and three $`p`$-like orbitals, each centered at the N atom. Two of the $`p`$-orbitals lie in the surface plane (the $`xy`$-plane) and are directed along the arrows shown in Fig. 13, while the third is oriented along the surface normal (the $`z`$-axis). We shall refer to the $`p`$-orbitals as $`p_x,p_y`$ and $`p_z`$, respectively. The energies corresponding to the atomic orbitals are (in eV): $`\epsilon _s=14.8`$, $`\epsilon _z=2.4`$, $`\epsilon _x=3.7`$, and $`\epsilon _y=4.2`$. We notice, that the energy of the $`p_x`$ and $`p_y`$ orbitals differ even though the symmetry of the system suggests that they should be equal. The reason for this is that the WFs break the four-fold rotation symmetry of the system, i.e. the subspace spanned by the four $`sp^3`$ WFs is not invariant under the same symmetry transformations as the Hamiltonian. This is not surprising, since the WFs are constructed solely from a criterion of maximal localization and no attempts are made to conserve symmetries. On the other hand we have found that by increasing the parameter $`E_0`$ above the value $`E_0=E_F`$ used in the present example, the symmetry between $`p_x`$ and $`p_y`$ can be restored. The price one has to pay is that the copper $`s`$-like WFs become less localized due to the further constrains on the localization space implied by the larger value of $`E_0`$. From the Hamiltonian in the WF basis we can also obtain the coupling, $`V`$, between each of the atomic nitrogen orbitals and its corresponding group orbital. These are quite similar and vary from $`3.1`$ eV to $`3.8`$ eV. In Fig. 14 we show the calculated PDOS for each of the three nitrogen $`p`$-orbitals (upper panel). Although the on-site energies of the $`p_x`$ and $`p_y`$ orbitals differ (as discussed above), their PDOS are rather similar. The PDOS of the corresponding group orbitals have been calculated with all coupling matrix elements to the N orbitals set to zero, i.e. the adsorbate states have effectively been decoupled from the surface. The result is shown in the lower panel of Fig. 14. For all three orbitals, the on-site energies lie within the band. Due to the strong coupling, bonding and anti-bonding resonances are formed at the band edges around $`7`$ eV and $`0`$ eV as can be seen in the upper panel of the figure. This is the limit of strong chemisorption. newns69 Since the four orbitals span all states with significant weight on the N atom, this representation provides a full representation of the nitrogen bonding. ## VI Conclusions We have presented a practical method for constructing partly occupied WFs for a wide range of systems. The method employs a bonding-antibonding closing procedure to filter out a set of unoccupied states, called the extra degrees of freedom, which serve to improve the localization of the WFs. The determination of the extra degrees of freedom is based on a minimization of the spread of the resulting WFs. We derived expressions for the gradients of the spread functional and showed how these can be combined with a Lagrange multiplier scheme to minimize the spread functional. The generality of the scheme was demonstrated by applying the method to three different systems. As an example of an isolated system, we considered a $`\text{Si}_5`$ cluster, and showed how different sets of WFs could be obtained by varying the number of extra degrees of freedom. A similar analysis was performed for a copper crystal, where we found results very similar to those of Souza et al. souza01 . Finally we studied in detail the WFs of a Cu(100) surface with a nitrogen coverage of 0.5. In many cases we were able to obtain a special set of WFs with a particularly high degree of symmetry and localization, by maximizing the average spread of the WFs. Moreover, the condition of maximal average localization was shown to coincide with a complete matching of bonding and antibonding states. ## VII Acknowledgments We acknowledge support from the Danish Center for Scientific Computing through Grant No. HDW-1101-05. ## Appendix A Spread functional for Vanderbilt ultrasoft pseudo-potentials For Vanderbilt ultra-soft pseudo-potentialsvanderbilt90 the optimal smoothness of the pseudo-wavefunctions is obtained by relaxing the norm-conserving constrains for the pseudo-wavefunctions. This results in a generalized orthonormality relationvanderbilt90 $$\psi _i\|S\|\psi _j=\delta _{ij}.$$ (36) The Hermitian operator $`S`$ is given by $$S=1+\underset{I}{}\underset{nm}{}q_{nm}\|\beta _n^I\beta _{m}^{}{}_{}{}^{I}\|,$$ (37) where the index $`I`$ denotes the atoms in the system, and $`q_{nm}`$ is given by $$q_{nm}=𝑑𝕣Q_{nm}^I(𝕣).$$ (38) The functions $`\{\beta _n^I\}`$ and $`\{Q_{nm}^I\}`$ are all localized functions centered at atom $`I`$. The functions $`\{Q_{nm}^I\}`$ describe the augmentation charge not contained in the smooth pseudo-wavefunctions, and they must therefore be included in the calculation of the spread of the wavefunctions. ### A.1 Large supercells In the case of large supercells, using the $`\mathrm{\Gamma }`$-point approximation, Bernasconi and Maddenbernasconi\_madden01 derived the following expression for the contribution to $`Z_\alpha ^{(0)}`$ from the augmentation charges $`Q_{nm}^I(𝕣)`$: $$Z_{\alpha ,ij}^{(us,0)}=\underset{I,nm}{}\psi _i\|\beta _m^I\beta _n^I\|\psi _j𝑑𝕣e^{\text{i}𝔾_\alpha 𝕣}Q_{mn}(𝕣)$$ (39) ### A.2 Periodic systems For the periodic case, using a uniform $`𝕜`$-point grid, we write $`Z_\alpha ^{us}`$ as $$Z_\alpha ^{us}=\underset{𝕜𝕜^{}}{}Z_\alpha ^{(us)𝕜𝕜^{}}.$$ (40) Here again we use the exact correspondence between the Bloch states $`\psi _{n𝕜}`$ and the $`\mathrm{\Gamma }`$-point eigenstates of the repeated cell. In the repeated cell we use the notation, $$h_{i𝕜}^{In𝕥}=\psi _{i𝕜}\|\beta _n^{I,𝕥}=\psi _{i𝕜}\|\beta _n^{I,𝕥=\mathrm{𝟘}}e^{i𝕜_𝕥}$$ (41) and $$Q_{nm}^{I𝕥}=\underset{𝔾}{}Q_{nm}(𝔾)e^{i𝔾(𝕣_𝕥)}$$ (42) $`_𝕥`$ is here a real space translation vector, given in terms of the basis $`𝕒`$, $`t_1𝕒_1+t_2𝕒_2+t_3𝕒_3`$, $`𝕥=(t_1,t_2,t_3)`$. We will use $`h^{In}=h^{I,𝕥=\mathrm{𝟘},n}`$ in what follows. Inserting $`h_{i𝕜}^{In𝕥}`$ and $`Q_{nm}^{I𝕥}`$ from Eqs. (41) and (42), together with the Bloch states, $`\psi _{i𝕜}`$, into the $`\mathrm{\Gamma }`$-point expression for $`Z_\alpha ^{(us,0)}`$ in Eq. (39), we find $$Z_{\alpha ,ij}^{(us,0),𝕜𝕜^{}}=\underset{𝕥}{}\underset{I,nm}{}h_{i𝕜}^{In𝕥}h_{j𝕜^{}}^{Im𝕥}\text{d}𝕣e^{i𝔾_\alpha 𝕣}Q_{nm}^{I𝕥}(𝕣)$$ (43) The sum over $`𝕥`$ is for $`t_i=0,..,N_i1`$, see Eq. (21). Inserting the left hand side of Eqs. (42) and (41), and rearranging, we find $$Z_{\alpha ,ij}^{(us,0),𝕜𝕜^{}}=\underset{I,nm}{}h_{i𝕜}^{In}h_{j𝕜^{}}^{Im}\underset{𝕥}{}e^{i(𝕜𝕜^{})_𝕥}\underset{𝔾}{}Q_{nm}(𝔾)e^{i𝔾_𝕥}\text{d}𝕣e^{i(𝔾𝔾_\alpha )𝕣}.$$ (44) Finally, we arrive at our expression for $`Z_{\alpha ,ij}^{(us,0),𝕜𝕜^{}}`$ $$Z_{\alpha ,ij}^{(us,0),𝕜𝕜^{}}=\underset{I,nm}{}h_{i𝕜}^{In}h_{j𝕜^{}}^{Im}\underset{𝕥}{}e^{i(𝕜𝕜^{}𝔾_\alpha )_𝕥}Q_{nm}(𝔾=𝔾_\alpha ),$$ (45) which is non-zero only when $`𝕜`$ and $`𝕜^{}`$ fulfills the condition in Eq. (29). Again we see that this expression contains the $`\mathrm{\Gamma }`$-point formalism, Eq. (39), as a special case.
warning/0506/cond-mat0506070.html	ar5iv	text	# Magnetic excitations in weakly coupled spin dimers and chains material Cu2Fe2Ge4O13 ## I introduction The ground states of low-dimensional magnets are strongly affected by quantum spin fluctuation. In so-called quantum spin liquids spin correlation remains short-range even at zero temperature, and the excitation spectrum is gapped. This disorder is relatively robust and such systems resist long-range ordering even in the presence of residual 3D interactions or anisotropy. In contrast, gapless low-dimensional magnets are very sensitive to external perturbations that can easily drive them towards long-range ordering. New physics is found in bicomponent systems that combine these two distinct types of low-dimensional spin networks. An example is $`R_2`$BaNiO<sub>5</sub> materials where Haldane spin chains weakly interact with magnetic rare earth ions. Zheludev96a ; Zheludev98a ; Alvarez et al. (2002) More recently, we reported the discovery and study of a novel quantum ferrimagnet Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub>.Masuda et al. (2003a, 2004) We showed that this compound can be viewed as a system of antiferromagnetic (AF) Cu-dimers that weakly interact with almost classical Fe<sup>3+</sup> chains. This weak coupling leads to a rather unusual cooperative ordering phase transition at low temperatures. The crystal structure of Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> is monoclinic $`P2_1/m`$ with $`a`$ = 12.101 Å, $`b`$ = 8.497 Å, $`c`$ = 4.869 Å, and $`\beta `$ = 96.131.Masuda et al. (2004) The arrangement of magnetic ions and likely exchange pathways is shown in Fig. 1 (a). Fe<sup>3+</sup> ions form crankshaft-shaped chains that run in the $`b`$ direction. These chains are separated by GeO<sub>4</sub> tetrahedra in $`c`$ direction. The Cu<sup>2+</sup> dimers are located in-between Fe<sup>3+</sup> chains along the $`a`$ direction. Simultaneous cooperative long-range magnetic ordering of the two magnetic subsystems occurs at 40 K.Masuda et al. (2004) The magnetic structure is roughly collinear, with spins lying in the crystallographic $`a`$ \- $`c`$ plane. In addition, there are small out-of-plane spin components (canting). The saturation magnetic moment on Cu<sup>2+</sup> ions is anomalously small, being suppressed by residual quantum fluctuations in the dimers: $`m_{\mathrm{Cu}}=0.38(4)\mu _\mathrm{B}`$. This suggests that the coupling $`J_{\mathrm{Cu}\mathrm{Fe}}`$ between the Cu-dimers to the Fe-subsystem is weak compared to intra-dimer AF interactions $`J_{\mathrm{Cu}}`$. The pairs of Cu-spins remain in a spin-singlet state that is only partially polarized by interactions with the long-range order in the Fe-subsystem. The data collected in preliminary inelastic neutron scattering experiments for energy transfers up to 10 meV could be well explained by fluctuations of Fe<sup>3+</sup> spins alone. The Fe<sup>3+</sup> spins form weakly-coupled $`S=5/2`$ chains, the corresponding effective exchange constants being $`J_{\mathrm{Fe}}=1.60(2)`$ meV, $`J_{\mathrm{Fe}}^{}=0.12(1)`$ meV, as shown in Fig. 1. These values are consistent with rough estimates of exchange constants based on magnetic susceptibility data: $`J_{\mathrm{Fe}}=1.7`$ meV, $`J_{\mathrm{Cu}}=25`$ meV. To date, no excitations associated with the Cu-dimers could be identified. While the very nature of the crankshaft-shaped chains implies some alternation of bond strength, we find no evidence thereof in the measured dispersion curves. In fact, they are well reproduced by a model involving uniform magnetic Fe-chains with weak interchain interactions along the $`c`$ direction. In the remainder of this work we shall therefore disregard the bond-alternation and assume the Fe-chains to be magnetically uniform. In the present paper we report a more detailed inelastic neutron scattering study of Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub>. Our main result is an observation of a complete separation of energy scales for Cu<sup>2+</sup>\- and Fe<sup>3+</sup>-centered magnetic excitations. An analysis of the intensity pattern for low-energy spin waves allows us to unambiguously associate them with the dynamics of $`S=5/2`$ chains. In addition, separate narrow-band excitation originating from Cu subsystem is observed at higher energy transfers. The layout of the paper will be as follows. In section II we will describe the characteristics of our samples and experimental setups. Section III describes the results for low-energy excitations in single-crystal samples, as well as measurements of the much weaker high energy excitations in a large-volume powder sample. In Section IV we shall interpret the observed separation of energy scales by a Mean Field- Random Phase Approximation (MF-RPA) treatment of interactions between Cu<sup>2+</sup>\- and Fe<sup>3+</sup> spins. Discussion and conclusion will be drawn in Section V and VI, respectively. ## II experimental High quality single crystals with the dimension of $`3\times 4\times 35`$ mm<sup>3</sup> were grown by the floating zone method. All samples were found to be twinned. The two twins share a common $`(b,c)`$ plane, the monoclinic structure allowing two independent orientations of the $`a`$ axis. In reciprocal space the domains share a common $`(a^{},b^{})`$ plane, but have distinct $`c^{}`$ axes, as illustrated in Fig. 1 (b). Two separate single crystals were co-aligned to obtain a larger sample of cumulative mosaic spread 0.44. For single crystal inelastic neutron scattering experiments we exploited four different setups. Setup I employed the SPINS cold neutron spectrometer at the NIST Center for Neutron Research (NCNR). The scattering plane was defined by the $`b^{}`$ and the bisector $`c^{{}_{}{}^{}}`$ of the $`c^{}`$ axes in the two crystallographic domains, as shown in Fig. 1 (b). In this geometry the scattering planes are the same in both domains. It is convenient to define $`c^{{}_{}{}^{}}c^{}\mathrm{cos}(\beta 90^{})`$. The momentum transfer in the scattering plane of the spectrometer is then indexed by Miller indexes $`h^{}`$, $`k^{}`$ and $`l^{}`$ of a fictitious orthorhombic structure. For the two domains $`h=\pm (c^{{}_{}{}^{}}/a^{})l^{}\mathrm{tan}(\beta 90^{}),k=k^{}`$ and $`l=l^{}`$. Since $`\beta `$ is close to 90, $`h`$ is almost zero for most measurements using Setup 1, where $`l`$ is small. In Setup 2 the scattering plane was $`(a^{},b^{})`$, which is also common for the two domain types. In both setups we used $`(\mathrm{guide})80^{}80^{}(\mathrm{open})`$ collimation with a Be filter positioned after the sample and a fixed final energy $`E_\mathrm{f}=5`$ meV or 3 meV. The data were collected at $`T`$ = 1.4 K using a standard He-flow cryostat. Setup 3 was used for wide surveys in reciprocal space and employed the HB1 thermal neutron spectrometer at the High Flux Isotope Reactor at ORNL. The scattering plane was $`(a^{},b^{}`$) , as in Setup 2. The collimation was $`48^{}40^{}40^{}240^{}`$. Neutrons with $`E_f`$ = 13.5 meV were used in conjunction with a Pyrolytic graphite (PG) filter positioned after the sample. The experiments in Setup 3 were performed at $`T=6.4`$ K maintained by a closed-cycle He refrigerator. Intensity was the main limiting factor for studies of higher-energy magnetic excitations. To maximize sample volume the measurements were performed on a 50 g polycrystalline Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> powder that was prepared by the solid state reaction method. This experiment was performed on the HB3 3-axis spectrometer at HFIR with $`48^{}40^{}60^{}120^{}`$ collimations and a PG filter after the sample (Setup 4). The final neutron energy was fixed at $`E_f`$ = 14.7 meV. A closed-cycle refrigerator was used to achieve low temperatures. ## III Experimental Result ### III.1 Low energies In this section we concentrate on the low energy excitations at energy transfers up to 10 meV. As was explained in our preliminary reportMasuda et al. (2004) and will be discussed in detail below, this part of the spectrum can be associated with conventional spin waves from the Fe-subsystem. The Cu-dimers only provide an effective Fe-Fe interaction, but contribute nothing to the dynamics at low energies. Typical energy scans at $`𝐪=(0k0.5)`$ measured using setup 1 are shown in Fig. 2. White and grey circles correspond to data collected with $`E_\mathrm{f}`$ = 5 meV and 3 meV respectively. Well defined resolution-limited peaks are observed in the entire Brillouin zone. Solid lines are Gaussian fits to the data after a subtraction of a linear background. The excitation energy is a minimum at the magnetic zone center at $`𝐪=(02.00.5)`$. A small gap of about 1.6 meV is observed at this wave vector. The apparent shoulder structure can be attributed to a splitting of the spin wave branch into two components with somewhat different gap energies. The gaps are most likely anisotropy-related. The zone-boundary energy at $`𝐪=(03.00.5)`$ is about 9 meV. The dispersion relation was obtained by Gaussian fits to individual scans and is shown in Fig. 3 (a). The $`k`$-dependence of the measured energy-integrated peak intensity is plotted in symbols in Fig. 3 (b). The intensity scales roughly as $`1/\omega `$. Energy scans at $`𝐪=(02.0l)`$ and $`(02.9l)`$ are shown in Fig. 4. Constraints on experimental geometry prevented us from reaching the more symmetric $`(03.0l)`$ reciprocal-space rods at higher energy transfers. These two sets of data reveal two distinct excitation branches, that are strongest near even and odd $`k`$-values, respectively. The former branch is dispersive along the $`c`$ axis, while the latter one is almost flat. Dispersion relations and the integrated intensity plots for both modes were obtained using Gaussian fits and are shown in Fig. 5 (a) and (b). Even for the more dispersive lower-energy branch the booundary energy along the $`c^{{}_{}{}^{}}`$ direction is only about 5 meV: about half of that in the $`b^{}`$ direction. No dispersion of low-energy excitations could be detected along the $`a^{}`$ direction. This is illustrated by the energy scans collected on the $`(h2.00)`$ reciprocal-space rod using Setup 2 (Fig. 6). Two peaks are observed at 5 and 9 meV, respectively. The excitation energies are plotted in Fig. 7 (a). Interestingly, the peak intensity is strongly dependent on $`h`$, as shown in Fig. 7 (b). From the data presented above one can immediately conclude that the most relevant magnetic interactions are within the Fe<sup>3+</sup> layers, parallel to the $`(b,c)`$ plane. Inter-layer coupling along the $`a`$-direction is considerably weaker. Based on structural considerations, it must involve the Cu<sup>2+</sup> spins. ### III.2 Higher-energy excitations The model that we previously proposed for Fe<sub>2</sub>Cu<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> (Ref. Masuda et al., 2004) implies a separation of energy scales of spin wave like excitations on the Fe subsystem and triplet excitations of the Cu dimers. From magnetic susceptibility measurements we have estimated the intra-dimer exchange constant to be around 25 meV. Due to small sample size, in single crystal experiments we failed to observe any clear magnetic inelastic features in this energy transfer range. To search for the Cu-triplet mode we performed additional measurements on a large-size powder sample. The powder data collected at $`T=12`$ K are summarized in the false color plot in Fig. 8 (a). The data were obtained by combining 17 separate constant-$`q`$ scans. For each such scan a linear background was subtracted from the measured intensity. Figure 8 (a) clearly shows a narrow excitation band at $`\mathrm{}\omega 24`$ meV. The measured intensity of the 24 meV peak decreases with the increase of $`q`$, as expected for magnetic scattering. The observed energy width at $`q=2.3`$ Å<sup>-1</sup> is somewhat larger than experimental resolution, but still small compared to the central energy (Fig. 8 (b)). Such energy dependence in powder samples typically indicates a narrow dispersion bandwidth (see, for example, Refs. Zheludev et al., 1996b; Masuda et al., 2003b). The observed spectrum is consistent with the excitations originating from the structural Cu-dimers, indicated by the thick solid bonds in Fig. 1 (a). However, the momentum transfer range covered in the experiment, especially at $`\|q\|0`$ is insufficient for a more detailed analysis of the structure factor. In particular, the size of the dimers can not be independently extracted from the experimental data. To obtain additional information on the 24 meV excitation we studied its temperature dependence at $`q`$ = 2.3 Å<sup>-1</sup>. The main challenge was dealing with a large background that originates from (i) temperature-independent scattering, including spurious scattering from an “accidental” Bragg powder line at 24.8 meV and (ii) temperature-dependent phonon scattering. These two contributions can be removed from the data if one assumes that the useful magnetic signal is relatively weak and/ or temperature-independent above $`T=200`$ K. This would indeed be true if the magnetic scattering originated from effectively isolated Cu-dimers with an intradimer exchange constant of 24 meV. The phonon contribution, which is assumed to scale with the Bose factor, is thus estimated from comparing scans at $`T=300`$ K and $`T=200`$ K. It is then appropriately scaled and subtracted from all scans, leaving only the true magnetic signal and the $`T`$-independent background. The two can not, in principle, be reliably separated. However, we can follow the change in magnetic signal from $`T=200`$ K by using the phonon-subtracted $`T=200`$ K scan as “background”. The result of this elaborate background subtraction is plotted in Fig. 9. The well-defined peak seen at low temperature broadens progressively with increasing $`T`$ and practically disappears at $`T40`$ K. In this range there also seems to be a downward shift in the peak’s central energy. ## IV Analysis ### IV.1 Separation of energy scales To understand the dynamics of the coupled Fe<sup>3+</sup> and Cu<sup>2+</sup> subsystems we shall make the central assumption that the exchange constant $`J_{\mathrm{Cu}}`$ that binds pairs of Cu<sup>2+</sup> spins into AF dimers is the largest energy scale in the system. Under these circumstances the degrees of freedom associated with Cu-spins can be effectively integrated out at low energies. Indeed, the isolated Cu<sup>2+</sup>-subsystem has a spin singlet ground state and a large energy gap $`\mathrm{\Delta }=J_{\mathrm{Cu}}`$. At $`\mathrm{}\omega \mathrm{\Delta }`$ it lacks any intrinsic dynamics, i.e., its dynamic susceptibility is purely real and almost energy-independent. In the spirit of RPA, from the point of view of Fe<sup>3+</sup> spins, the Cu-dimers merely act as a polarizable medium that can transfer magnetic interactions between the Fe-layers. The staggered spin susceptibility of each $`S=1/2`$ dimer being $`2/(J_{\mathrm{Cu}})`$, this effective coupling, labeled as $`J_{\mathrm{eff}}`$ in Fig. 1 is given by: $$J^{\mathrm{eff}}=J_{\mathrm{Cu}\mathrm{Fe}}^2/(2J_{\mathrm{Cu}}).$$ (1) Thus, to a good approximation, the low-energy spin dynamics of Fe<sub>2</sub>Cu<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> is simply that of the Fe<sup>3+</sup>-subsystem with an additional exchange coupling. From the experiment, where no dispersion of spin waves could be observed along the $`a`$-axis, the effective exchange constant must be rather small. Nevertheless, as explained in our previous paper,Masuda et al. (2004) it is absolutely crucial in completing a 3-dimensional spin network and allowing long-range magnetic ordering at a non-zero temperature. At high energy transfers, comparable to the Cu-dimer gap, it is the Fe<sup>3+</sup> degrees of freedom that can be effectively integrated out. At $`T<T_\mathrm{N}`$ their effect is reduced to producing a static staggered exchange spin field that acts on the Cu<sup>2+</sup> spins and is proportional to the ordered Fe<sup>3+</sup> moment: $`h^{\mathrm{Cu}}`$ $`=`$ $`S_{\mathrm{Fe}}J_{\mathrm{Cu}\mathrm{Fe}}.`$ (2) In our approximation at high energy transfers Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> behaves as a collection of (possibly interacting) Cu-dimers in an effective staggered field.Masuda et al. (2004) ### IV.2 Spin waves and dynamic structure factor for low energies As explained above, the low-energy spectrum of Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> can be understood by considering the Fe-subsystem in isolation, and even $`J_{\mathrm{eff}}`$ can be ignored for its smallness. For the magnetically ordered state the dynamic structure factor of such a system can be calculated using conventional spin wave theory (SWT). For simplicity we shall assume a collinear magnetic structure, ignoring the small canting of the Fe<sup>3+</sup> spins out of the $`(a,c)`$ plane. Even though there are four Fe ions in each unit cell of Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub>, the topology of exchange interactions defined by $`J_{\mathrm{Fe}}`$ and $`J_{\mathrm{Fe}}^{}`$ in Fig. 1 is that of a regular Bravais lattice. This equivalency is illustrated in Fig. 10. The exact relation between the actual structure factor $`S(𝐪,\omega )`$ and $`S_0(𝐪,\omega )`$ for the equivalent lattice is readily obtained: $`S(𝐪,\omega )`$ $`=`$ $`S_0(𝐪,\omega )\mathrm{cos}^22\pi \delta _ah\mathrm{cos}^22\pi \delta _bk+{\displaystyle \frac{1}{2}}S_0(𝐪+(0,1,0),\omega )(1+\mathrm{sin}^22\pi \delta _ah\mathrm{sin}^22\pi \delta _bk)`$ (3a) $`+`$ $`S_0(𝐪+(0,2,0),\omega )\mathrm{cos}^22\pi \delta _ah\mathrm{sin}^22\pi \delta _bk+{\displaystyle \frac{1}{2}}S_0(𝐪+(0,3,0),\omega )(1+\mathrm{sin}^22\pi \delta _ah\mathrm{sin}^22\pi \delta _bk),`$ $`\delta _a=0.1239,\delta _b=0.0629.`$ (3b) Here $`\delta _a`$ and $`\delta _b`$ are the displacements of Fe<sup>3+</sup> ions from high-symmetry positions, as shown in Fig. 1 (b). The SWT dynamic structure factor for a collinear antiferromagnet on a Bravais lattice is well known:Lovesey $`S_0(𝐪,\omega )`$ $`=`$ $`(u_𝐪+v_𝐪)^2\delta (\mathrm{}\omega _𝐪\mathrm{}\omega )`$ (4a) $`u_𝐪^2`$ $`=`$ $`{\displaystyle \frac{S(\mathrm{}\omega _𝐪+2Sj(0))}{\mathrm{}\omega _𝐪}},`$ (4b) $`u_𝐪v_𝐪`$ $`=`$ $`{\displaystyle \frac{2S^2j(𝐪)}{\mathrm{}\omega _𝐪}},`$ (4c) $`\mathrm{}\omega _q`$ $`=`$ $`S\sqrt{j(0)^2j(𝐪)^2+\mathrm{\Delta }^2}.`$ (4d) Here $`j(𝐪)`$ is the Fourier transform of exchange interactions, and $`\mathrm{\Delta }`$ empirically accounts for the anisotropy gap. In our particular case of the Fe<sup>3+</sup> subsystem in Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> we have: $$j(𝐪)=2(J_{\mathrm{Fe}}\mathrm{cos}\frac{\pi }{2}k+J_{\mathrm{Fe}}^{}\mathrm{cos}2\pi l),$$ (5) The cross section for inelastic neutron scattering from spin waves is given by: $$\frac{d^2\sigma }{d\mathrm{\Omega }dE}\|F(𝐪)\|^2\left[1+\left(\frac{q_y}{q}\right)^2\right]n_𝐪+1S(𝐪,\omega ),$$ (6) In this formula $`F(𝐪)`$ is the magnetic form factor for Fe<sup>3+</sup>, $`n_𝐪`$ is the Bose factor and $`q_z`$ is the projection of the scattering vector onto the direction of ordered Fe<sup>3+</sup> moments. ### IV.3 Fits to data The model cross section given by Eq. (6) can accurately reproduce the observed low-energy spectra in Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub>. Due to the presence of four terms in Eq. (3), there are four distinct spin wave branches, that we shall denote as modes I through IV, correspondingly. The dispersion relation given by Eq. (4d) was fit to the experimental data shown in Figs. 3 (a) and 5 (a) using a least-squares algorithm. A good fit is obtained with $`J_{\mathrm{Fe}}=1.60`$ meV, $`J_{\mathrm{Fe}}^{}=0.12`$ meV, and $`\mathrm{\Delta }=2.02`$ meV. The result is shown in lines in figures 3a and 5a. With these parameters our model also agrees well with the measured dispersion (or, rather, absence thereof) along the $`a`$ axis, as shown in 7a. What is important, is that not only the energies, but also the intensities of the observed excitations are well reproduced by our model. Calculated intensities for each mode or combined intensities of a couple of modes in cases where experimental energy resolution is insufficient to resolve individual branches, are shown in lines in Figs. 3b, 5b and 7b. This quantitative agreement between the measured and observed structure factors confirms that the scattering is indeed due to Fe<sup>3+</sup> spins. An excellent illustration of this was obtained by mapping out the scattering intensity in a wide $`𝐪`$-range using Setup 4, as shown in Fig. 11 (a). These data correspond to a fixed energy transfer $`\mathrm{}\omega =8.5`$ meV, and have a characteristic checkerboard pattern. Our spin wave model with the parameters quoted above reproduces this behavior very well, as shown in Fig. 11 (b). In this calculation the model cross section was numerically convoluted with the known experimental resolution function. The apparent periodicity along $`k`$ is due to a steep dispersion that takes the excitations in and out of the probed energy range. However, the periodicity along $`h`$ is related to the trigonometric coefficients in Eq. (3a). These, in turn, are determined by the geometry of the crankshaft-shaped Fe<sup>3+</sup> chains in Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub>. ## V Discussion As demonstrated above, the low-energy spin dynamics of Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> is well described by an effective spin wave theory for the Fe<sup>3+</sup> spin chains. Based on the available data it is impossible to unambiguously associate the observed 24 meV mode with the Cu<sup>2+</sup> dimers. However, much confidence in this assumption can be drawn from a recent study of Cu<sub>2</sub>Sc<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub>.Masuda05unpublished In this isostructural compoundRedhammer and Roth (2004) only the Cu<sup>2+</sup> ions are magnetic. Indeed the magnetic susceptibility showed $`S=1/2`$-dimers behavior with similar energy to Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub>.Masuda05unpublished We can now estimate all the relevant exchange interactions in the system. Using Eq. (2) in combination with the staggered susceptibility of an isolated antiferromagnetic dimer, from the known saturation moments $`S_{\mathrm{Cu}}=0.18(8)`$ and $`S_{\mathrm{Fe}}=1.77(4)`$ for Cu<sup>2+</sup> and Fe<sup>3+</sup>, respectively, Masuda et al. (2004) we get $`J_{\mathrm{Cu}\mathrm{Fe}}=2.54(3)`$ meV. The effective coupling is then $`J^{\mathrm{eff}}=0.13(4)`$ meV. The results for all exchange parameters are summarized in Table 1. Our model for Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> is qualitatively consistent with the observed slight increase of the energy of the Cu-dimer mode with decreasing temperature. Indeed, the gap energy of isolated dimers, as that of other gapped systems, is known to increase with the application of a staggered field.Jolicoeur and Golinelli (1994); Ma et al. (1995) In Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub> this field is generated by the ordered moment on the Fe<sup>3+</sup> sites, and at $`T<T_\mathrm{N}`$ increases proportionately to $`S_{\mathrm{Fe}}`$. Beyond that, the observed $`T`$-dependence of the 24 meV mode is different from that for isolated dimers. In the latter, the intensity would remain almost constant below $`T=40`$ K. Moreover, the peak would remain sharp at all temperatures. The discrepancy may be related to intrinsic limitations of the MF/RPA approach, and merits further investigation. In particular, it is tempting to somehow associate the observed emergence and sharpening of the 24 meV inelastic peak with the onset of long-range order. ## VI Conclusion Our data bring solid quantitative support to the concept of separation of energy scales in the mixed-spin quantum antiferromagnet Cu<sub>2</sub>Fe<sub>2</sub>Ge<sub>4</sub>O<sub>13</sub>. Using a simple MF/RPA approach we are able to determine all the relevant exchange interactions. However, certain features of the temperature dependence of spin excitations require further theoretical and experimental study. ###### Acknowledgements. Work at ORNL was carried out under Contracts No. DE-AC05-00OR22725, US Department of Energy. Experiments at NIST were supported by the NSF through DMR-0086210 and DMR-9986442.
warning/0506/cond-mat0506491.html	ar5iv	text	# Knudsen Effect in a Nonequilibrium Gas ## Abstract From the molecular dynamics simulation of a system of hard-core disks in which an equilibrium cell is connected with a nonequilibrium cell, it is confirmed that the pressure difference between two cells depends on the direction of the heat flux. From the boundary layer analysis, the velocity distribution function in the boundary layer is obtained. The agreement between the theoretical result and the numerical result is fairly good. molecular dynamics simulation, nonequilibrium steady state, heat conduction, boundary layer analysis, Boltzmann equation, Knudsen effect Although there has been a long history of nonequilibrium statistical mechanics since Boltzmann introduced the Boltzmann equation, the understanding of nonequilibrium statistical mechanics is still in the primitive stages . The significant role of nonequilibrium physics in the mesoscopic region has been recently recognized. For example, there has been some important progress such as the Fluctuation Theorem and the Jarzynski equality in mesoscopic nonequilibrium statistical mechanics. In typical situations of mesoscopic physics, materials are confined to narrow regions. Thus the boundary effects at the mesoscopic scale are important not only for nonequilibrium statistical mechanics but also for consideration of friction and lubrication . On the other hand, some macroscopic phenomenologies for nonequilibrium steady states have been proposed. These are the Extended Thermodynamics (ET) , the Extended Irreversible Thermodynamics (EIT) which is the combination of ET and information theory, and the Steady State Thermodynamics (SST) . It is interesting that both EIT and SST treat a common process in which a nonequilibrium cell is connected with an equilibrium cell . The Knudsen effect is the phenomenon in which two equilibrium cells with different temperatures are connected by a small hole . The balance equation of the Knudsen effect is given by $$\frac{P_1}{\sqrt{T_1}}=\frac{P_2}{\sqrt{T_2}},$$ (1) where $`T_i`$ and $`P_i`$ represent the temperature and the pressure in the cell, $`i`$, respectively. Although eq. (1) contains only macroscopic variables, it is in contrast to the ordinary thermodynamic balance condition where the pressures the two cells are equal. Such an exceptional condition means that the mass balance is determined by the transportation of the gases in the small hole. Therefore, the relevance of predictions by macroscopic theory for the generalized Knudsen effect in which an equilibrium cell is connected with a nonequilibrium cell by a small hole is questionable. On the other hand, the explicit perturbative solution of the Boltzmann equation for hard spheres has been derived at the Burnett order of the heat flux . The quantitative accuracy and numerical stability of their solution in the bulk region has been confirmed by molecular dynamics (MD) simulation for hard-spheres and the extension to the tenth order shows that their second order solution is accurate even when the heat current is large . They have confirmed that the solution in the bulk region derived by information theory is not consistent with that of the Boltzmann equation . Kim and Hayakawa have also discussed the nonequilibrium Knudsen effect. Their result denies the prediction of SST, but both theories predict that the osmosis $`\mathrm{}P`$ defined by the $`\mathrm{}PP_{xx}^{neq}P^{eq}`$ with the $`xx`$ component of the pressure tensor $`P_{xx}^{neq}`$ in the nonequilibrium cell and the pressure $`P^{eq}`$ in the equilibrium cell is always positive regardless of the direction of heat flux. However, their simplification using the bulk solution of the Boltzmann equation in the boundary layer is not acceptable. In fact, numerical simulations of the Boltzmann equation of hard spheres exhibit the discontinuity of velocity distribution function (VDF) in the boundary layer, but there is no discontinuity of VDF derived by Kim and Hayakawa . Furthermore, it is also well known that the gas temperature near the wall is different from the wall temperature , but both treatments ignore this fact. To clarify the truth of the nonequilibrium Knudsen effect, we employ the event-driven molecular dynamics simulation of hard-disks developed in refs. \[20-23\]. Let $`x`$ and $`y`$ be the Cartesian coordinates of the horizontal and the vertical directions, respectively. First we check the validity of eq. (1) by MD. We adopt the diffusive reflection for the vertical walls away from the hole, the simple reflection rule for the wall between two cells and the periodic boundary condition for the horizontal wall. We connect two cells by a small hole, as illustrated in Fig. 1. For $`T_1/T_2=1.96`$, we have found the steady value of $`P_1/P_2\sqrt{T_2/T_1}=0.9982\pm 0.0195`$, where the average area fraction and the number of the particle are $`0.015`$ and $`10,000`$ respectively. We determine the pressure based on the Virial theorem and average in $`10^4`$ collisions per particle. The stationary state is realized after $`3\times 10^5`$ collisions per particle and we use the data after that. We simulate the system until $`10^6`$ collisions have been performed per particle. Next, we simulate the nonequilibrium Knudsen effect. The number of hard disks and the average area fraction are the same values used to simulate the conventional Knudsen effect. We adopt the diffusive boundary condition for the wall between two cells (see Fig. 1). We examine the values of $`T_1/T_2`$ as $`1.96`$ and $`1/1.96`$. We cause both cells to divide into 10 equal parts, and we have confirmed that pressure based on the Virial theorem is identical in each cell. This is consistent with the stationary condition . In nonequilibrium cases the stationary state is realized after $`10^6`$ collisions per particle and we use the data after that. We simulate the system until $`4.4\times 10^6`$ collisions have been performed per particle. Our results of MD plotted in Fig. 2 indicate that the sign of $`\mathrm{}P(P_{xx}^{neq}P^{eq})`$ depends on the direction of energy flux. Actually, the stationary value $`\mathrm{}P`$ in our simulation is given by $`\mathrm{}P/<n_2T_2>=0.0279\pm 0.0184`$ for $`T_1/T_2=1.96`$, and $`\mathrm{}P/<n_2T_2>=0.0152\pm 0.0175`$ for $`T_1/T_2=1/1.96`$, where $`<n_2T_2>`$ is the time average of $`n_2T_2`$ (Fig. 2). This result contrasts with the positive $`\mathrm{}P`$ predicted by both SST and Kim and Hayakawa . Now let us compare the VDF of MD with the perturbative solution of the Boltzmann equation at the Burnett order obtained by Kim for 2D hard disks. From now on, we restrict our interest to the data for heat flux $`J<0`$, with $`T_1/T_2=1/1.96`$. As shown in Fig. 3, VDF obtained from MD in the bulk region of a nonequilibrium cell is almost identical to the theoretical VDF, where the VDF of MD is obtained from the average of particles existing in $`\pm 1.05d`$ from the center of nonequilibrium cell in the horizontal direction with hard disk diameter $`d`$. On the other hand, the VDF in the boundary layer in which we average the data of particles existing between $`0.7d`$ and $`2.8d`$ apart from the right wall of the nonequilibrium cell deviates from the theoretical VDF , particularly for $`v_x<0`$ (Fig. 4). Let us derive VDF in the boundary layer in the nonequilibrium cell. We assume that the VDF of the incident particles ($`v_x>0`$) in the boundary layer is the same as the bulk distribution function ($`f_{NE}`$). On the other hand, we assume that the VDF for the particles reflected by the wall ($`v_x<0`$) obey the Maxwellian. Thus, the distribution function $`f_{BL}`$ in the boundary layer is given by $$f_{BL}(𝐯)=\{\begin{array}{cc}f_{MB}(n_w,T_w,𝐯)\hfill & v_x<0\hfill \\ f_{NE}(n_x,T_x,J,𝐯)\hfill & v_x0,\hfill \end{array}$$ (2) where $$f_{MB}(n_w,T_w,𝐯)\frac{n_wm}{2\pi T_w}\mathrm{exp}\left[\frac{m𝐯^2}{2T_w}\right],$$ (3) with the mass of a hard-disk $`m`$, and $$\begin{array}{cc}\hfill f_{NE}(n_x,T_x,J,& 𝐯)f_{MB}(n_x,T_x,𝐯)\times \hfill \\ & \left[1\frac{mJv_x}{2b_1n_xT_x^2}\underset{r1}{}r!b_rL_r^1\left(\frac{m𝐯^2}{2T_x}\right)\right],\hfill \end{array}$$ (4) with $`b_1=1.03,b_2=5.738\times 10^2,b_3=4.946\times 10^3,b_4=4.313\times 10^4,b_5=3.452\times 10^5,b_6=2.241\times 10^6`$ and Laguerre’s bi-polynomial $`L_b^a(x)`$ . Since the heat flux is sufficiently small, we adopt the first order nonequilibrium VDF in the heat-flux for $`f_{NE}`$. There are three unknown variables, $`n_w`$, $`n_x`$, and $`T_x`$ in eqs. (3) and (4), while there are three relations, $`n`$ $`{\displaystyle 𝑑𝐯f_{KL}(𝐯)},`$ (5) $`J`$ $`m{\displaystyle 𝑑𝐯v_x\frac{𝐯^2}{2}f_{KL}(𝐯)},`$ (6) $`{\displaystyle 𝑑𝐯v_xf_{KL}(𝐯)}=0.`$ (7) Here, the first two equations are definitions of the density $`n`$ and the heat flux $`J`$, and the last equation represents the mass balance condition. Therefore, we can determine $`n_w`$, $`n_x`$, and $`T_x`$ from eqs. (5)-(7). The expansions in terms of $`J`$ of the three variables become $`n_w=n+a{\displaystyle \frac{m^{1/2}}{T_w^{3/2}}}J,`$ (8) $`n_x=n+b{\displaystyle \frac{m^{1/2}}{T_w^{3/2}}}J,`$ (9) $`T_x=T_w+c{\displaystyle \frac{m^{1/2}}{nT_w^{1/2}}}J,`$ (10) where $`a`$, $`b`$ and $`c`$ are constants to be determined. From eqs. (5),(6) and (7), we obtain the solutions of the linear simultaneous equations as $`(a,b,c)(0.32,0.099,0.84)`$. The distribution function $`f_{BL}`$ near the wall has thereby been determined by $`n,T_w`$, and $`J`$. Figure 4 is the comparison of $`f_{BL}`$ with the results of MD. From Fig. 4, both $`f_{BL}`$ and the VDF from MD have a discontinuity at $`v_x=0`$, as in the case of conventional boundary layer analysis . For reflective VDF ($`v_x<0`$), there is the small difference between the result of MD and the Maxwellian. We may deduce that it arises from collisions between particles because we measure the VDF a short distance from the wall. With the aid of $`f_{BL}`$, $`\mathrm{}P`$ can be calculated as $$\mathrm{}P\frac{a+b+c+0.412}{2}\frac{J}{\sqrt{T_2}}=0.415\frac{J}{\sqrt{T_2}}.$$ (11) From this equation and the value of the heat-flux in MD, we evaluate $`\mathrm{}P/<n_2T_2>0.0298`$ for $`T_1/T_2=1.96`$ and $`\mathrm{}P/<n_2T_2>0.0156`$ for $`T_1/T_2=1/1.96`$. The result agrees well that of the simulation. We can also obtain the temperature jump coefficient $`\gamma 0.97`$ which is defined through $`T_{x=0}T_w=\gamma J/n\sqrt{m/T_w}`$. Now, let us discuss our result. There are some advantages to employing MD as the numerical simulation. First, we can easily change boundary conditions for walls depending on our interest. The system of connecting two cells by a small hole is easily simulated by MD. Second, MD is suitable for high density simulation of gases. The nonequilibrium Knudsen effect for high density gases is an interesting subject for future discussions. We are also interested in the size effect of the hole on the transition from Knudsen balance to the thermodynamic balance. On the other hand, there are some disadvantages to employing MD. Because of the small system size of our MD simulation, the fluctuation of the pressure is large. The advantage of our boundary layer analysis is that we can write the explicit form of VDF in the boundary layer in terms of the density, the heat-flux and the temperature of the wall. On the other hand, VDF by Sone et al has an implicit form that is obtained as a numerical solution of an integral equation. The explicit VDF can be obtained by our simplification, but the validity of this method has not been confirmed. In fact, for 3D hard-sphere gases, our method predicts the temperature jump coefficient $`\gamma 0.72`$, while Sone et al predicts $`\gamma 1.00`$. To check the validity of our boundary layer analysis, it will be necessary to employ 3D MD simulation for 3D. We also stress the difficulty of describing the nonequilibrium Knudsen effect by macroscopic phenomenology. In fact, our boundary layer analysis strongly depends on the boundary condition. In conclusion, the sign of $`\mathrm{}P`$ depends on the direction of the heat flux. The approximated VDF in the boundary layer has been obtained from the assumption that the incident particles obey the bulk VDF and the reflected particles obey the Maxwellian. This agrees well with the simulation result. The authors would like to thank M. Fushiki and H.-D. Kim for their valuable comments. They appreciate Aiguo Xu for his critical reading of the manuscript. This work is partially supported by a Grant-in-Aid from the Japan Space Forum and the Ministry of Education, Culture, Sports, Science and Technology (MEXT), Japan (Grant No. 15540393) and a Grant-in-Aid from the $`21`$st century COE “Center for Diversity and Universality in Physics” from MEXT, Japan.
warning/0506/cond-mat0506405.html	ar5iv	text	# Symbiotic Solitons in Heteronuclear Multicomponent Bose-Einstein condensates ## I Introduction Symbiosis is an assemblage of distinct organisms living together. Although the original definition of symbiosis by De Bary 1879 did not include a judgment on whether the partners benefit or harm each other, currently, most people use the term symbiosis to describe interactions from which both partners benefit. In Physics, waves in dispersive linear media tend to expand due to the different velocities at which the wave components propagate. This is not the case in many nonlinear media, in which certain wavepackets, called solitons are able to propagate undistorted due to the balance between dispersion and nonlinearity Scott . Stable solitons of different subsystems are sometimes able to “live together” and form stable complexes called vector solitons as it happens with Manakov optical solitons Manakov ; Kivshar or stabilized vector solitons PGVector . In some cases, a (large) robust soliton can be used to stabilize a (small) weakly unstable wave saturable . Multicomponent solitary waves also appear in Bose-Einstein condensates (BECs). In fact, multicomponent BECs support nonlinear waves which do not exist in single component BECs such as domain wall solitons Cohen ; Kasamatsu , dark-bright solitons Anglin , etc. Most of the previous analyses correspond to homonuclear multicomponent condensates for which the atom-atom interactions are repulsive. However, heteronuclear condensates offer a wider range of possibilities, the main one being the possibility of having a negative inter-species scattering length. This possibility has been theoretically explored in the context of Feschbach resonance management Simoni and realized experimentally for boson-fermion mixtures KRb ; NaLi . In this paper we study the existence and properties of bright solitons in heteronuclear two-component BECs with scattering lengths $`a_{11},a_{22}>0`$ and $`a_{12}<0`$. We would like to stress the fact that these coefficient combinations do not arise in other systems where similar model equations are used. For instance in nonlinear optics, where the nonlinear Schrödinger equations used to describe the propagation of laser beams in nonlinear media are similar to the mean field equations used to describe Bose-Einstein condensates, the nonlinear coefficients are allways of the same sign. The closest analogy could happen in the so-called QPM ( quasi-phase-matched ) quadratically nonlinear media, where an effective cubic nonlinearity could be “engineered” which could have similar properties but we do not know of any systematic studies of those systems. Our analysis will show novel features with respect to those already found in single species BECs bright . For instance, even when solitons do not exist for each of the species, the coupling leads to robust vector solitons. Since the mutual cooperation between these structures is essential for their existence we will refer to these solitons hereafter as symbiotic solitons. We also show how they appear by modulational instability and study some features of their collisions. We also comment on the possibility of obtaining these structures in multidimensional configurations. ## II The model and its basic properties In this paper we will study two-component BECs in the limit of strong transverse confinement ruled by Perez-Garcia $`i{\displaystyle \frac{u_1}{t}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2u_1}{x^2}}+\left(g_{11}\|u_1\|^2+g_{12}\|u_2\|^2\right)u_1,`$ (1a) $`i{\displaystyle \frac{u_2}{t}}`$ $`=`$ $`{\displaystyle \frac{\alpha }{2}}{\displaystyle \frac{^2u_2}{x^2}}+\left(g_{21}\|u_1\|^2+g_{22}\|u_2\|^2\right)u_2,`$ (1b) where $`x`$ is the adimensional longitudinal spatial variable measured in units of $`a_0=\sqrt{\mathrm{}/m_1\omega _{}}`$, $`t`$ is the time measured in terms of $`1/\omega _{}`$, and $`u_j(𝐱,t)u_j(𝐫,\tau )\sqrt{a_0^3}`$. The dimensional reduction leads to Perez-Garcia $`g_{ij}=2a_{ij}\alpha ^{i+j2}/a_0`$, with $`\alpha =m_1/m_2`$ and $`a_{ij}`$ being the $`s`$wave scattering lengths. The normalization for $`u_j`$ is $`\|u_j\|^2d^3x=N_j`$ where $`N_j`$ is the number of particles of each species. Let us first consider constant amplitude solutions of Eq. (1), which are of the form $`\varphi _j(z,t)`$ $`=`$ $`A_je^{i\beta _jt},`$ (2a) $`\beta _j`$ $`=`$ $`g_{jj}\|A_j\|^2+g_{j,3j}\|A_{3j}\|^2,`$ (2b) for $`j=1,2`$. We will study the evolution of small perturbations of $`\varphi _j`$ of the form $$u_j(z,t)=\left(A_j+\delta A_j(z,t)\right)e^{i(\beta _jt+\delta \beta _j(z,t))}$$ (3) Using Eq. (1) and retaining the first order terms we get partial differential equations for $`\delta A_1,\delta \beta _1,\delta A_2,\delta \beta _2`$ which can be transformed to Fourier space to obtain $$\delta A_j(z,t)=_{}a_0(k)e^{ikx}e^{\mathrm{\Omega }(k)t}𝑑k$$ (4) $`a_0(k)`$ being the Fourier transform of the initial perturbation. Perturbations remain bounded if $`\text{Re}\left[\mathrm{\Omega }(k)\right]0`$. Some algebra leads to $$\mathrm{\Omega }^2=\frac{1}{2}\left(f_1+f_2\pm \sqrt{(f_1f_2)^2+4C^2}\right)$$ (5) where $`f_j=(g_{jj}A_j^2+k^2/4)k^2,C^2=A_1^2A_2^2g_{12}^2k^4`$. The so-called modulational instability (MI) occurs when $`\mathrm{\Omega }(k)^2>0`$ for any $`k`$. For small wavenumbers (worst situation) we get $$g_{12}^2>g_{11}g_{22},$$ (6) which is analogous to the miscibility criterion for two-component condensates Kasamatsu . However, the physical meaning of Eq. (6) is very different since now this instability is a signature of the tendency to form coupled objects between both atomic species. The role of MI in the formation of soliton trains and domains in BEC has been recognized in previous papers Kasamatsu ; Min1 ; bright . ## III Vector solitons Eqs. (1) have sech-type solutions $$u_j(x,t)=\left(\frac{N_j}{2\omega }\right)^{1/2}\text{sech}\left(\frac{x}{\omega }\right)e^{i\lambda _jt}$$ (7) with $`\lambda _1=1/(2\omega ^2),\lambda _2=\alpha /(2\omega ^2)`$, and $`\omega =2/\left(g_{11}N_1g_{12}N_2\right),`$ provided the restriction $$g_{12}\left(m_1N_2m_2N_1\right)=m_2g_{22}N_2m_1g_{11}N_1,$$ (8) and the MI condition (6) are satisfied. Eq. (8) implies that, given the number of particles in one component the other is fixed. Since the self-interaction coefficients are positive, these solitons are supported only by the mutual attractive interaction between both components. This type of vector soliton thus differs from others described for Nonlinear Schrödinger equations of the form Eq. (1), such as the Manakov solitons Manakov , where all the nonlinear coefficients cooperate to form the solitonic solution. The MI condition (6) implies that the formation of these solitons has a threshold in $`g_{12}`$ and means that the cross-interaction must be strong enough to be able to overcome the self-repulsion of each atomic cloud. There are no analogues to this condition in single component systems since solitons exist for any value of the self-interaction coefficient $`g<0`$. To fix ideas, taking a <sup>87</sup>Rb-<sup>41</sup>K mixture with $`a_{11}=69a_0,a_{22}=99a_0`$ the MI condition implies that $`a_{12}<83a_0`$ in order to obtain solitons. In Fig. 1(a) it can be seen how the ratio $`N_2/N_1`$ is close to 0.4 in the range of values of $`83a_0>a_{12}>150a_0`$. An hypothetical <sup>7</sup>Li-<sup>23</sup>Na mixture with $`a_{11}=5a_0`$ and $`a_{22}=52a_0`$ (in appropriate quantum states) leads to the curve in Fig. 1(b), which shows a much larger range of variation. ## IV Soliton Stability We can use the Vakhitov-Kolokov (VK) criterion to study the stability of solitons given by Eq. (7). To do this, we must study the sign of $`\lambda _j/N_j`$. For soliton solutions this can be done from the explicit form of $`\lambda _j`$. After some algebra we find $`\lambda _1(N_1)`$ and $`\lambda _2(N_2)`$ and obtain that $`\lambda _1/N_1>0,`$ and $`\lambda _2/N_2>0`$ in all their range of existence, which proves the linear stability of the solitons for small perturbations and contradicts the naive intuition that the self-repulsion would lead to intrinsically unstable wavepackets. We have studied numerically the robustness of symbiotic solitons to finite amplitude perturbations. First we have perturbed both solutions with small amplitude noise and found that, in agreement with the predictions of the VK criterion, they survive after the emission of the noise in the form of radiation. Next we have applied a stronger perturbation consisting of displacing mutually their centers and observe that a soliton is formed even for relative displacements of the order of the soliton size \[Fig. 2\]. Finally we have started with sech-type initial data which are not solitons and observe that after the emission of some radiation solitons are formed. ## V Generation of symbiotic solitons by MI To study the generation of these solitons by MI in realistic systems we have considered a multicomponent Bose-Einstein condensate of <sup>87</sup>Rb and <sup>41</sup>K atoms for which the inter-species scattering length $`a_{12}`$ is controlled by the use of Feschbach resonances as proposed in Simoni . To simplify the problem here we do not consider the effect of gravity. We start by constructing the ground state of the system for an elongated trap typical of the LENS setup with $`\omega _{}=`$ 215 Hz, $`\omega `$ = 16.3 Hz. For these atomic species $`a_{11}=69a_0`$ and $`a_{22}=99a_0`$. We adjust the inter-species scattering length to $`a_{12}=95a_0`$ during the condensation process. The ground state of this system for $`N_1=25000,N_2=20000`$, shown in Fig. 3(a), agrees well with the theoretical predictions for these systems BT . After the condensate is formed we change instantaneously this quantity to a negative value and at the same time switch off the longitudinal trapping potential and observe numerically the evolution of the ground state. First we choose $`a_{12}=90a_0`$ and observe the evolution starting from the ground state with $`a_{12}=95`$. Since the inter-component repulsive force is not present now, the sharp domain wall separating both species (see Fig. 3(a)) decay through a highly oscillatory process related to the formation of a shock wave Kuzmiak . The final outcome is the formation of a soliton train (see Fig. 3(b,c)) of which three solitons of about 20 $`\mu `$m size and each with about 3000 rubidium and 1200 potassium atoms remain in our simulation domain after 500 adimensional time units \[Fig. 3(c)\]. Other smaller and wider solitons exit our integration region traveling at a faster speed. The final number of solitons depends on the value of $`a_{12}`$ choosen during the condensation process (which controls the overlapping of the species) and the number of particles $`N_1`$, $`N_2`$ and the negative scattering length $`a_{12}`$ choosen to destabilize the system. For instance, choosing $`a_{12}=70a_0`$, which is below the theoretical limit for MI the evolution of the wavepacket is purely dispersive \[see Fig. 4(a)\]. Choosing $`a_{12}=87a_0`$, above the MI limit but below the choice of Fig. 3 leads to the formation of a single soliton \[Fig. 4(b)\]. It seems that the larger the scattering length, the larger the number of solitons which arise after the decay of the initial configuration. The many degrees of freedom present in these system open many posibilities for controling the number and sizes of solitons by appropriately choosing the values of $`a_{12}`$ before and after the condensate is released and the initial number of particles $`N_1,N_2`$. ## VI Collisions of symbiotic solitons The robustness of symbiotic solitons manifests also in their collisional behavior and their internal structure makes the interaction of these vector solitons very rich. Since each soliton is a compound object the collisions are at the same time a coherent phenomenon because of the direct overlapping of the same type of atoms and an incoherent one because of the incoherent nature of interaction between different types of atoms. A related subject of recent interest in Optics is that of partially coherent solitons partial . We have simulated head-on collisions of equal symbiotic solitons of opposite velocities given by $`u_j`$ $`=`$ $`\sqrt{{\displaystyle \frac{N_j}{2w}}}\text{sech}\left({\displaystyle \frac{x+x_0}{w}}\right)e^{iv\sqrt{m_j}x+i\alpha _{j,+}}`$ (9) $`+`$ $`\sqrt{{\displaystyle \frac{N_j}{2w}}}\text{sech}\left({\displaystyle \frac{xx_0}{w}}\right)e^{iv\sqrt{m_j}x+i\alpha _{j,}}`$ for $`j=1,2`$. $`\alpha _{j,\pm }`$ are the relative phases and $`N_2`$ is given by Eq. (8). In Fig. 5 we show some examples of these collisions. Slow \[Fig. 5(a-d)\] or moderate speed collisions \[Fig. 5(e-f)\] lead to bound solitons while for larger speeds the picture is not so clear. The specific outcome of the collision depends on the relative soliton phases with the phase difference between the larger components in the symbiotic soliton (in this case Rb) being the dominant ones. For instance collisions with phases $`𝜽(\alpha _{1,+},\alpha _{1,},\alpha _{2,+},\alpha _{2,})=(0,\pi ,\pi ,0)`$ \[Fig. 5 (c)\] and $`𝜽=(0,0,\pi ,0)`$ (not shown) both lead to mutual repulsion but the outgoing speeds are different due to the different interactions between the internal components of the soliton. Collisions with higher but still moderate speeds \[Fig. 5(e-f)\] give independent vector solitons. The outcome of the collisions with zero phases is a bound state of two vector solitons which has internal oscillations, i.e. some sort of multicomponent higher order soliton. ## VII Prospects for Multidimensional Symbiotic Solitons A very interesting question arising naturally is: do these symbiotic solitons exist in multidimensional scenarios? In principle the answer is not evident since the only effect acting against stabilization of multidimensional soliton structures would be collapse, but one could think that in this case collapse could be inhibited because of the repulsive self-interaction, thus a deeper analysis is in order. The adimensional model equations in two and three dimensions take the form $$i\frac{u_j}{t}=\left(\frac{1}{2m_j}\mathrm{\Delta }+V_j+g_{j,j}\|u_j\|^2+g_{j,k}\|u_k\|^2\right)u_j,$$ (10) with $`j=1,2`$ and $`k=2,1`$ correspondingly. Let us first consider this problem in two spatial dimensions. To study collapse rigorously one usually tries to compute the exact evolution of the wavepacket widths rigorously PG99 . For the multicomponent case and $`m_1=m_2=m`$, this was studied by group-theoretical methods by Gosh . In our case, from the general formulae obtained by Gosh we get a sufficient condition for collapse, which is $$\begin{array}{c}=_^n[\underset{j=1,2}{}(\|u_j\|^2/(2m)+V_j\|u_j\|^2\hfill \\ \hfill +g_{jj}\|u_j\|^4/2)+g_{12}\|u_1\|^2\|u_2\|^2]<0.\end{array}$$ (11) In principle, this is a bad result for obtaining localized structures since it means that arbitrarily close to any stationary solution (for which $`=0`$), there would be collapsing solutions and thus stationary solutions, if they exist, would be unstable. As it is usual in the framework of collapse problems the situation would be even worse in three spatial dimensions with solutions of arbitrary small number of particles undergoing collapse provided they are initially sufficiently localized. This means that in principle symbiotic solitons could only be obtained in quasi-1D geometries because of the transverse stabilization effect provided by the trap in a similar way as ordinary bright solitons do. ## VIII Conclusions and extensions In this paper we have studied vector solitons in heteronuclear two-component BECs which are supported by their attractive mutual interaction. These symbiotic solitons are linearly stable and remarkably robust and can be generated through modulational instability phenomenon with many possibilities for control. Collisions of these vector solitons show their robustness and open different ways for their manipulation and the design of novel quantum states such as breather-like states. We have also considered multidimensional configurations and shown that collapse may avoid the formation of fully multidimensional symbiotic solitons. We think that the conceptual ideas behind our work can also be used to understand boson-fermion mixtures. For instance, $`a_{12}`$ is known to be negative and large for quantum degenerate mixtures of <sup>87</sup>Rb and <sup>40</sup>K FB2 . In those systems numerical simulations have proven the formation of localized wavepackets Karpiuk which could share the same essential mechanisms for the formation of solitary waves. ###### Acknowledgements. We acknowledge V. Vekslerchik, R. Hulet and B. Malomed for discussions.This work has been partially supported by grant BFM2003-02832 (Ministerio de Educación y Ciencia, Spain).
warning/0506/cond-mat0506398.html	ar5iv	text	# Multi-Scaling of Correlation Functions in Single Species Reaction-Diffusion Systems ## I Introduction Modeling chemical reactions is an important practical and theoretical problem. Systems of reacting particles are typical of complex irreversible non-equilibrium systems. The popular description of these systems in terms of simple rate equations currently adopted in chemical kinetics house does not appear to have a firm theoretic foundation and often produces wrong results. Interacting particle systems provide a good model for simple chemical reactions. Well-known examples include systems of diffusing-coalescing and diffusing-annihilating particles describing reactions $`A+AA`$ and $`A+A\mathrm{}`$ respectively. Extensive studies of these particle systems in low dimensions have shown that the rate equations yield incorrect results for the computation of average concentrations of reactants, see priv for review. Rate equations fail in low dimensions due to the presence of large fluctuation effects, which violate mean-field theory (MFT) assumptions underlying their derivation. The study of diffusive annihilation (or coalescence)<sup>1</sup><sup>1</sup>1Both systems belong to the same universality class Cardy in a sense that the correlation functions for both systems are identical apart from the amplitude. For more details see mass2 . is a good starting point for analyzing large fluctuation effects in more complicated non-equilibrium statistical systems. For example, the system $`A+AA`$ was used to analyze the aggregation of massive diffusing particles Oleg2 . The aim of the current paper is to study the effects of large fluctuations on correlation functions of an arbitrary order for the reactions $`A+A\mathrm{}`$ in $`d2`$ and $`A+A+A\mathrm{}`$ in $`d=1`$. Large fluctuation effects are accounted for in binary reaction-diffusion models using Empty Interval methods (EIM) and its generalizations mass2 ; mass1 ; ben ; Doer . This approach is restricted to $`d=1`$ and does not extend to higher dimensions. There are rigorous results on the average density of particles in $`d=1,2`$ Bram . The Smoluchowski approximation gives correct answers for average concentrations Howd but cannot be used for higher order correlation functions. In the early 90’s, the work of Cardy and Lee Cardy ; Lee used field theoretic methods, in particular the renormalization group (RG) to obtain an answer for the average density as well as its amplitude for $`d2`$. The study of $`A+A\mathrm{}(A)`$ has also introduced the concept of stochastic rate equations with imaginary multiplicative noise Cardy . In this paper we derive the multi-scaling of correlation functions in the systems $`A+A\mathrm{}`$ and $`3A\mathrm{}`$ using the RG method. The main object of our study is the large time temporal scaling of the probability of finding $`N`$ particles in a small volume $`\mathrm{\Delta }V`$, denoted $`P_t(N,\mathrm{\Delta }V)`$. For the $`A+A\mathrm{}`$ reaction-diffusion system we are interested in $`d2`$. We do not consider higher dimensions as the answers there are given by MFT. Most studies have concentrated on computing the average density of particles ($`N=1`$)Cardy ; Lee . To the best of our knowledge, the computation of multi-particle probabilities are only considered in mass2 ; mass1 ; ben , with the analysis restricted to one dimension. Dynamical RG method allows us to obtain answers for the large time limit in the form of an $`\epsilon `$-expansion ($`\epsilon =2d`$) for $`d<2`$ and logarithmic corrections to the MF scaling for $`d=2`$: $$\frac{P_t(N,\mathrm{\Delta }V)}{P_t(1,\mathrm{\Delta }V)^N}\{\begin{array}{cc}t^{\frac{N(N1)\epsilon }{4}+𝒪(\epsilon ^2)}\hfill & d<2\hfill \\ (\mathrm{ln}t)^{\frac{N(N1)}{2}}\left(1+𝒪(\frac{1}{\mathrm{ln}t})\right)\hfill & d=2\hfill \end{array}$$ (1) where Lee $$P_t(1,\mathrm{\Delta }V)\{\begin{array}{cc}t^{d/2}\hfill & d<2\hfill \\ \frac{\mathrm{ln}t}{t}\hfill & d=2\hfill \end{array}$$ Equation (1) represents the multi-scaling or the deviation of $`P_t(N,\mathrm{\Delta }V)`$ from $`P_t(1,\mathrm{\Delta }V)^N`$. As $$\underset{t\mathrm{}}{lim}\frac{P_t(N,\mathrm{\Delta }V)}{P_t(1,\mathrm{\Delta }V)^N}=0$$ (2) equation (1) reflects the anti-correlation between particles in the large time limit. For the ternary reaction-diffusion system $`3A\mathrm{}`$ we restrict our attention to $`d=1`$. MFT provides the answers in higher dimensions Lee . The average density has been studied in Lee ; Krap , and the two-point function in Lee . Higher order correlations have not been considered. The method of empty intervals does not extend to the ternary reaction. The RG method produces asymptotically exact results for $`d=1`$: $$\frac{P_t(N,\mathrm{\Delta }V)}{P_t(1,\mathrm{\Delta }V)^N}\mathrm{ln}t^{\frac{N(N1)(N2)}{6}}$$ (3) where Lee $$P_t(1,\mathrm{\Delta }V)\left[\frac{\mathrm{ln}t}{t}\right]^{1/2}$$ The paper is organized as follows: an introduction to the lattice model of $`A+A\mathrm{}`$ is given in Section II. A, its field-theoretic re-formulation is given in Section II.B, followed by a summary of the mean field results in Section II.C. Section III contains the RG analysis of the binary annihilation system for $`d2`$, starting with the description of renormalization procedure in Section III.A. Section III.B contains the derivation of the corresponding Callan-Symanzik equation for the theory, its solution and the derivation of large time asymptotics of multi-particle probabilities. We prove the exactness of our result for $`P_t(N=2,\mathrm{\Delta }V)`$ using the first Hopf equation of the theory. In Section III.C we compare our results in $`d=1`$ against results from mass2 and establish their equivalence for $`N=1,2,3,4`$. Further we conjecture the exactness of our one-loop answer in $`d=1`$ for general values of $`N`$. In Section IV we extend our analysis to the ternary reaction in $`d=1`$ using the same methodology. ## II Field-Theoretic Formulation and Mean-Field limit of $`A+A\mathrm{}`$ model ### II.1 The Model Consider a set of point particles performing random walks, characterized by diffusion coefficient D, on the lattice $`𝐙^d`$. Any two particles positioned at the same site can annihilate each other according to an exponential process with rate $`\lambda `$. For the simplicity of our analysis we assume finite reaction rates. However the large time asymptotics of our model belongs to the universality class of instantaneous annihilation-diffusion model (see end of Section III B). It is assumed that the initial distribution of particle number at each site is independent Poisson with mean $`N_o`$. Let the random variable $`N_t(x)`$ represent the occupation number for site $`x`$ at time $`t`$. The configuration vector $`\underset{¯}{N}\{N(x)\}_{x𝐙^d}`$ specifies the state of the system at time $`t`$ by encoding the occupation number at all sites. $`\underset{¯}{N}`$ is also termed a microstate. Let $`𝒫_t(\underset{¯}{N})`$ be the probability of finding the system in microstate $`\underset{¯}{N}`$ at time $`t`$. Correlation functions of $`N_t(x)`$ can be obtained by averaging functions of $`N(x)`$ with respect to $`𝒫_t(\underset{¯}{N})`$. For example, the average density is given by: $$\overline{N_t(x)}=\underset{[\underset{¯}{N}]}{}N(x)𝒫_t(\underset{¯}{N})$$ (4) Due to translational invariance this will be independent of $`x`$. Our main object of interest is the large time limit asymptotics of $`𝒫_t(N(x)=N)`$, the probability of finding $`N`$ particles at site $`x`$. It turns out that in the low density limit this probability is proportional to the $`N^{th}`$ factorial moment of $`N_t(x)`$, which we denote by $`M_N(x,t)`$. This can be verified as follows: $`M_N(x,t)`$ $`=`$ $`\overline{[N_t(x)(N1)]\mathrm{}[N_t(x)1]N_t(x)}`$ (5) $`=`$ $`{\displaystyle \underset{[\underset{¯}{N}]}{}}{\displaystyle \underset{k=0}{\overset{N1}{}}}[N(x)k]𝒫_t[\underset{¯}{N}]`$ $`=`$ $`{\displaystyle \underset{[\underset{¯}{N}]}{}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{N1}{}}}[nk]\chi _{[N(x)=n]}𝒫_t[\underset{¯}{N}]`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{N1}{}}}[nk]𝒫_t[N(x)=n]`$ $`=`$ $`{\displaystyle \underset{n=N}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{N1}{}}}[nk]𝒫_t[N(x)=n]`$ $``$ $`N!𝒫_t[N(x)=N]`$ In the above derivation, the first five lines are exact relations. The last line is due to the fact that in the large time limit we expect the particle density to be low and particles are anti-correlated Cardy ; mass2 ; ben ; Lee . Hence the configurations with the smallest possible value of $`N_t=N`$ will give the dominant contribution. The coarse-grained counterpart of $`𝒫_t[N(x)=N]`$ is $`P_t(N,\mathrm{\Delta }V)`$, the probability of finding $`N`$ particles in the volume element $`\mathrm{\Delta }V`$. Let $`\mathrm{\Delta }N_t(x)`$ be the number of particles in a volume $`\mathrm{\Delta }V`$ (centred at $`x`$) at time t: $$\mathrm{\Delta }N_t(x)=_{\mathrm{\Delta }V}d^dyn_t(y)$$ (6) where $`n_t(y)`$ stands for the $`density`$ of particles at time $`t`$ at a point $`y`$. It should be mentioned that upon the averaging $`\overline{\mathrm{\Delta }N_t}`$ is independent of $`x`$ due to translational invariance. In the limit of large time and fixed $`\mathrm{\Delta }V`$, factorial moments of $`\mathrm{\Delta }N`$ are related to $`P_t(N,\mathrm{\Delta }V)`$ via a relation analogous to (5): $`P_t(N,\mathrm{\Delta }V)`$ $`=`$ $`{\displaystyle \frac{1}{N!}}M_N(t)`$ (7) $`=`$ $`{\displaystyle \frac{1}{N!}}\overline{{\displaystyle \underset{k=0}{\overset{N1}{}}}[\mathrm{\Delta }N_t(x)k]}`$ As we will demonstrate in the next section, the factorial moments of $`\mathrm{\Delta }N_t(x)`$ admit a simple representation in terms of polynomial moments of Doi’s fields. ### II.2 Path-Integral representation In the last section it was shown that to obtain correlation functions of interest, we need to know moments of $`N(x)`$ with respect to $`𝒫_t(\underset{¯}{N})`$. It is possible to evaluate these moments in a path-integral setting. The formalism is due to Doi-Zeldovich-Ovchinnikov and we refer the reader to Cardy for details. We simply present a schematic derivation. The time evolution of the probability measure $`𝒫_t`$ on the space of microstates is given by the master equation. The master equation is a linear first order autonomous differential equation with respect to time, which implies that its evolution operator can be expressed as a path-integral. Thus we can find a path-integral representation for any correlation function. As we are interested in universal properties of the system at scales much larger than the lattice spacing, it is convenient to work with the continuum (coarse-grained) limit of the lattice model, which corresponds to some effective field theory in imaginary time. The continuum limit in the path-integral is taken according to the rules $$\lambda (\mathrm{\Delta }x)^d\lambda D(\mathrm{\Delta }x)^2D$$ where $`\mathrm{\Delta }x`$ denotes the lattice spacing. The reader is asked to refer to Lee for more details. The resulting field theory is given by the effective action $`S`$: $$S=_0^\tau 𝑑td^dx\overline{\varphi }(_t\mathrm{\Delta })\varphi +2\lambda \overline{\varphi }\varphi ^2+\lambda \overline{\varphi }^2\varphi ^2n_o\overline{\varphi }\delta (\tau )$$ (8) where $`n_o`$ is the initial average particle density. We will be working in the large $`n_o`$ limit. The diffusion coefficient D can be set to 1 as it is not renormalized by fluctuations Lee . The $`\overline{\varphi }`$-field is a response field. The $`\varphi `$-field is related to the local density field, $`n_t(x)`$ in the sense that there is a one-to-one correspondence between correlation functions of $`n_t`$ and $`\varphi `$. Let $``$ be the averaging with respect to the path-integral measure $`e^S`$. Then the average density is Lee $$\overline{n_t(x)}=\varphi (x,t)=𝒟\overline{\varphi }(x,t)𝒟\varphi (x,t)\varphi (x,t)e^{S[\varphi ,\overline{\varphi }]}$$ (9) Let us introduce the quantity $`\mathrm{\Delta }\varphi `$: $$\mathrm{\Delta }\varphi (x)=_{\mathrm{\Delta }V}d^dy\varphi (y)$$ (10) where the volume $`\mathrm{\Delta }V`$ is centred at $`x`$. Averaging (10) results in an $`x`$-independent quantity due to the translational invariance of the system. Let us consider the statistics of $`\overline{\mathrm{\Delta }N_t(x)}`$ defined in equation (6). Moments of $`\mathrm{\Delta }N`$ and moments of $`\mathrm{\Delta }\varphi `$ are related as follows $`\mathrm{\Delta }\varphi `$ $`=`$ $`\overline{\mathrm{\Delta }N_t}`$ (11) $`\mathrm{\Delta }\varphi ^2`$ $`=`$ $`\overline{\mathrm{\Delta }N_t(\mathrm{\Delta }N_t1)}`$ (12) The relations (11), (12) follow from Lee . Let us note that (11) gives us a path-integral representation of the average number of particles in a volume $`\mathrm{\Delta }V`$. More generally Oleg1 $$\mathrm{\Delta }\varphi ^N=M_N(t)$$ (13) where $`M_N`$ is defined by equation (7). In the previous section it was shown that in the large-time limit the factorial moment is approximately equal to the probability of finding $`N`$ particles in a volume element $`\mathrm{\Delta }V`$. Hence, $$P_t(N,\mathrm{\Delta }V)=\frac{1}{N!}\mathrm{\Delta }\varphi ^N$$ (14) As we mentioned before, we are interested in the limit $`\mathrm{\Delta }V0`$. Physically this limit corresponds to $`\mathrm{\Delta }Vl^d`$, where the correlation length $`l\sqrt{t}`$ Cardy . By working in the large time limit we are effectively dealing with the case $`\mathrm{\Delta }V0`$. Then we may write: $$\mathrm{\Delta }\varphi ^N=_{\mathrm{\Delta }V}𝑑x_1\mathrm{}𝑑x_N\varphi (x+x_1)\mathrm{}\varphi (x+x_N)$$ (15) where $`\varphi (x+x_1)\mathrm{}\varphi (x+x_N)`$ $`=`$ $`𝐅(x_1\mathrm{}x_N)\varphi ^N(x)_R`$ (16) $`(1+𝒪(\mathrm{\Delta }V^{1/d}))`$ $`𝐅(x_1\mathrm{}x_N)`$ is the leading coefficient of Wilson’s operator product expansion (OPE) Jean . The limit $`\mathrm{\Delta }V0`$ corresponds to the ultraviolet ($`x_i^{}s0`$) asymptotics of $`\varphi (x+x_1)\mathrm{}\varphi (x+x_N)`$, which is given by the renormalized average of the composite operator $`\varphi ^N`$. We conclude that multi-scaling of $`P_t(N,\mathrm{\Delta }V)`$ is described by anomalous dimensions of the composite operators $`\varphi ^N`$. The spectrum of anomalous dimensions of $`\varphi ^N`$ will be computed in Section III. ### II.3 Mean Field Limit and Loop expansion The Feynman rules for the perturbative computation of correlation functions are derived from the action $`S`$, equation (8). They are given by Fig 1. The propagator $`G_o`$ is the Green’s function of the standard diffusion equation. The perturbative expansion of $`\varphi ^m`$ in powers of $`\lambda `$ is given by the sum of all diagrams with $`m`$ outgoing lines; hence diagrams contributing to the mean density, $`\varphi `$, have one outgoing line. The action $`S`$ must be dimensionless; thus in length (L) units, the dimensions of the relevant parameters must be $$[t]=L^2[\overline{\varphi }]=L^0[\varphi ]=L^d[\lambda ]=L^{d2}.$$ (17) The critical dimension is $`d_c=2`$, where the reaction rate is dimensionless. For $`d<2`$ the field theory is super-renormalizable where all Feynman integrals converge. For $`d=2`$ we use dimensional regularization to ensure convergence of Feynman integrals. The bare dimensionless reaction rate is given by $`g_o(t)=\lambda t^{1d/2}`$, which grows with time in $`d<2`$. A combinatorial argument shows that an $`n`$-loop diagram contributing to the mean density is proportional to $`g_o^{n1}(t)`$ Lee ; thus in the weak coupling regime the main contribution to the mean density comes from the sum of tree diagrams. This is equivalent to the mean-field approximation. This sum (tree diagrams) is also termed the classical density, denoted $`n_{cl}(t)`$. In $`d<2`$ MFT is valid for small times given initial density is large. This agrees with our intuition: for small times local fluctuations around the large mean value of the density are small. At large times, $`g_o(t)`$ grows and MFT breaks down. To compute corrections to the mean-field answers, one must compute higher loop contributions. In particular this involves summing infinite sets of diagrams for a fixed number of loops. These can be re-summed in a more compact form using the classical density ($`n_{cl}`$) and the classical response function ($`G_{cl}`$). The classical response function consists of the sum of all tree diagrams with one outgoing and one incoming line. The diagrammatic form of the integral equations satisfied by these two quantities are given in Fig 2. The solution to these equations are: $`n_{cl}(t)`$ $`=`$ $`{\displaystyle \frac{1}{2\lambda t}}`$ (18) $`G_{cl}(x_2,t_2;x_1,t_1)`$ $`=`$ $`\left[{\displaystyle \frac{n_{cl}(t_2)}{n_{cl}(t_1)}}\right]^2G_o(x_2,t_2;x_1,t_1)`$ (19) Using $`n_{cl}`$ and $`G_{cl}`$ in combination with vertices of Fig 1, we arrive at Feynman rules generating finitely many diagrams for a given number of loops. We use $`n_{cl}`$ to write a mean-field expansion for $`\varphi ^m`$ which is valid for small times in $`d2`$. Using dimensional analysis it can be verified that a diagram with $`m`$ outgoing lines ($`\varphi ^m`$) and $`n`$ loops is proportional to $`n_{cl}^mg_o^n`$. Then, $$\varphi ^m=n_{cl}^m(t)\left(1+\underset{n=1}{\overset{\mathrm{}}{}}c_{m,n}g_o^n(t)\right)$$ (20) In this form we see that for small $`g_o(t)`$ the loop corrections are small. Then we may formulate the following mean-field answer for the probability of finding $`N`$ particles in volume $`\mathrm{\Delta }V`$: $`P_t(N,\mathrm{\Delta }V)`$ $`=`$ $`{\displaystyle \frac{1}{N!}}\mathrm{\Delta }\varphi ^N`$ (21) $`\stackrel{MFT}{=}`$ $`{\displaystyle \frac{1}{N!}}\mathrm{\Delta }\varphi ^N`$ $``$ $`(\mathrm{\Delta }V)^Nt^N`$ Comparing (21) with results from mass2 ; mass1 ; ben in $`d=1`$ we find the linear temporal scaling of MFT to be incorrect in the large time limit. We will compute the correct scaling in the subsequent Section. ## III Renormalization Group Analysis of $`P_t(N,\mathrm{\Delta }V)`$ The dynamical renormalization group method allows one to extract large time asymptotics of correlation functions of the theory (8). The first step is to eliminate all $`\epsilon 0`$ singularities of Feynman integrals at some reference time $`t_o`$. The process of removing these divergences is called renormalization. This is done by introducing a renormalized reaction rate, renormalized fields, etc. The number of renormalization constants needed to eliminate all divergences is finite for $`dd_c`$ due to renormalizability of (8). Individual terms in the renormalized perturbative series for any correlation function $`C(t)`$ depend on the reference time $`t_o`$. The lack of dependence of the unrenormalized version of $`C(t)`$ on the unphysical parameter $`t_o`$ leads to renormalization group (Callan-Symanzik) equation for the correlation function. This is solved subject to the initial condition at $`t_o`$ given by the perturbative expansion for $`C(t)`$ at $`t=t_o`$. This procedure is equivalent to re-summing all leading $`\epsilon `$-singularities in $`d<d_c`$ at all orders of the loop expansion Coll . Consequently, one obtains scaling laws for correlation functions in terms of an $`\epsilon `$-expansion. The knowledge of $`𝒪(\epsilon )`$ terms in the loop expansion in $`d<d_c`$ yields leading order logarithmic corrections to the mean field answers in $`d=d_c`$. We will now apply the method described above to compute both temporal and spatial scaling exponents of $`P_t(N,\mathrm{\Delta }V)`$ for $`A+A\mathrm{}(d_c=2)`$ reaction. ### III.1 One-loop renormalization of composite operator $`\varphi ^N`$ Dimensional analysis shows that renormalization of the correlation function $`_{i=1}^N\varphi (x_i,t)`$, where $`x_ix_j`$, requires reaction rate renormalization only Cardy ; Lee . However, we are interested in single-point correlation functions of the form $`\varphi ^N(x,t)`$, see (16). The operator $`\varphi ^N`$ is called a composite operator for $`N2`$. It is well known that insertion of composite operators under the sign of averaging leads to new types of divergences in the corresponding loop expansion Jean . These divergences cannot be eliminated by reaction rate renormalization and require multiplicative renormalization of the corresponding composite operators. As we will see below, these extra divergences are responsible for the multi-scaling of probabilities $`P_t(N,\mathrm{\Delta }V)`$. For the theory (8) the first instance of such a divergence occurs in the correlator $`\varphi ^2(y,t)_c=lim_{xy}\varphi (x,t)\varphi (y,t)_c`$ as $`\epsilon 0`$. We will see that it is specifically this singularity that leads to the multi-scaling of $`P_t(N,\mathrm{\Delta }V)`$ for $`N2`$. The tree-level answer is quoted below for $`d<2`$Lee ; Mun : $$\varphi (x,t)\varphi (y,t)_c=\frac{1}{8\pi g_ot^d}\left[\frac{1}{\epsilon }\left[\left[\frac{\|xy\|}{\sqrt{t}}\right]^\epsilon 1\right]+O(\epsilon ^0)\right]$$ (22) The averaging $`\mathrm{}_c`$ denotes connected correlation functions. It is clear from equation (22) that for $`xy`$ there is no divergence in the limit $`\epsilon 0`$. However if $`x=y`$, there will be a $`\frac{1}{\epsilon }`$ divergence (as $`\epsilon 0`$) despite (22) being a tree-level answer. This divergence cannot be regularized by reaction rate renormalization. It requires multiplicative renormalization, where the divergent function is multiplied by a renormalizing factor which will eliminate the singularity. This procedure violates naive dimensional arguments according to which $`\varphi ^N`$ should scale as $`\varphi t^{Nd/2}`$ for $`d<2`$. Multiplicative renormalization leads to a non-trivial anomalous dimension which depends non-linearly on the order of the composite operator. Multiplicative renormalization is not required for the average density $`\varphi `$ which explains its lack of anomalous scaling Cardy . Let $`M_N(t)=\varphi ^N(t)`$. Our aim is to renormalize $`M_N`$ at (arbitrary reference) time $`t_o`$. Let $`g_o=\lambda t_o^{\epsilon /2}`$ be the bare (dimensionless) reaction rate. The renormalized rate $`g`$ is also defined at reference time $`t_o`$ and is given by Lee : $$g=\frac{g_o}{1+g_o/g^{}}$$ (23) where $`g^{}=\frac{\epsilon }{C_d}=2\pi \epsilon +𝒪(\epsilon ^2)`$ is the non-trivial stable fixed point of the RG flow in the space of effective reaction rates. The constant $`C_d`$ is given by $$C_d=\frac{2\epsilon }{(8\pi )^{d/2}}\mathrm{\Gamma }\left(\frac{\epsilon }{2}\right)$$ (24) $`C_d`$ is regular at $`d=2`$ and takes the value $`\frac{1}{2\pi }`$. Expressing $`g_o`$ as a power series in $`g`$ and substituting into $`M_N(t_o)`$, enables us to obtain $`M_N`$ as a power series in $`g`$: $$M_N(t_o)=\left[n_{cl}(g_o,t_o)\|_{g_o=g}\right]^N\left[1+\underset{n=1}{\overset{\mathrm{}}{}}c_{N,n}g^n\right]$$ (25) The expression for $`M_N(t_o)`$ to $`𝒪(g^{1N})`$ is summarized by the diagrams in Fig. 3. The first term in this series is just $`n_{cl}^N`$ and is given by Fig 3a. Expanding $`n_{cl}`$ in $`g`$ using relation (23) will generate a term proportional to $`\frac{1}{g^{}}\frac{1}{\epsilon }`$. This term acts as a counter term to the diagram shown in Fig 3b. It eliminates the singularity in the one-loop correction to mean density $`n^{(1)}`$ Lee . The first divergence which cannot be eliminated by the renormalized reaction rate, appears in the coefficient $`c_{N,1}`$. The origin of this singularity is from averaging the composite operator $`\varphi ^2`$. Let $`Z_N`$ be the renormalizing factor of $`M_N`$. $$Z_N=1+\underset{n=1}{\overset{\mathrm{}}{}}a_{N,n}g^n$$ (26) It is chosen in such a way that the renormalized correlation function $`M_{N,R}`$ is non-singular at $`\epsilon =0`$ at all orders. By definition <sup>2</sup><sup>2</sup>2Lack of operator mixing produces this simple form for equation (27). This is due to the fact that in our theory all processes (diagrams) reduce the particle number. For more see Oleg1 . $$M_{N,R}(g,t_o)=Z_N(g,\epsilon )M(g,t_o,\epsilon )$$ (27) The coefficients $`a_{N,n}`$ are chosen to cancel the singularities in the coefficients $`c_{N,n}`$. This is known as the minimal subtraction scheme. In particular $$a_{N,1}=\left[c_{N,1}\right]_s$$ (28) where $`[]_s`$ extracts the singular part of an expression. Let us look at the calculation in more detail. Using Fig. (3) we may write $$M_N(t_o)=\left[\frac{1}{2gt_o^{d/2}}\right]^N\left[1\frac{N(N1)}{4\pi \epsilon }g+𝒪(g^2)+\mathrm{finite}\right]$$ (29) The $`𝒪(g)`$ term is computed as follows: we need $`n_{cl}^{(N2)}`$, there are $`{}_{N}{}^{}C_{2}^{}`$ diagrams of type c and finally by setting $`x=y`$ in (22) which gives $`[M_2(t_o)]_s=\frac{1}{8\pi \epsilon }t_o^d`$. Hence $$Z_N=1+\frac{N(N1)}{4\pi \epsilon }g+𝒪(g^2)$$ (30) ### III.2 RG computation of Multi-Scaling The knowledge of renormalization laws (23) and (30) can be used to compute $`M_{N,R}(t)`$ for $`t>t_o`$ as follows. The Markov property of the evolution operator $`U`$ for the bare theory, namely $`U(t,t_o)U(t_o,0)=U(t,0)`$, tells us that the bare function $`M_N(t)`$ is independent of $`t_o`$ for $`t>t_o`$. Hence $$t_o\frac{}{t_o}M(t)=t_o\frac{}{t_o}\left[Z_N^1M_{N,R}(t,t_o)\right]=0$$ (31) The function $`M_{N,R}(t)`$ is a function of $`(t,t_o,g)`$, leading to the ansatz: $`M_{N,R}(t)=t_o^{Nd/2}\mathrm{\Phi }(\frac{t}{t_o},g(t_o))`$, where $`\mathrm{\Phi }`$ is a function with dimensionless arguments. The choice of $`t_o`$ is arbitrary, but we choose it small enough so that MFT is still valid. This motivates the choice of the pre-factor $`t_o^{Nd/2}`$. Upon substitution of this ansatz into (31) we obtain the Callan-Symanzik equation: $$\left[t\frac{}{t}+\beta (g)\frac{}{g}+\frac{Nd}{2}+\gamma _N(g)\right]M_{N,R}(t,t_o,g)=0$$ (32) where the $`\beta `$ and $`\gamma _N`$ functions of the theory are given by $`\beta (g)`$ $`=`$ $`t_o{\displaystyle \frac{g}{t_o}}={\displaystyle \frac{1}{2}}\left[C_dg^2\epsilon g\right]`$ (33) $`\gamma _N(g)`$ $`=`$ $`\beta (g){\displaystyle \frac{\mathrm{ln}Z_N}{g}}={\displaystyle \frac{N(N1)}{8\pi }}g+𝒪(g^2)`$ (34) The initial condition (at time $`t_o`$) is given by the loop expansion of $`\varphi ^N`$ with the most dominant contribution coming from the MFT answer: $$M_{N,R}(t_o,g)=n_{cl}^N(t_o,g)$$ (35) Equation (32) subject to initial condition (35) is solved using the method of characteristics and has the following solution for $`d<2`$: $`M_{N,R}(t,g(t,t_o))`$ $`=`$ $`\left({\displaystyle \frac{t_o}{t}}\right)^{Nd/2}n^N(t_o,g(t,t_o))`$ (36) $`\left[{\displaystyle \frac{g(t,t_o)g^{}}{gg^{}}}\right]^{\frac{N(N1)}{4\pi \epsilon }g^{}}`$ $`g(t,t_o)`$ $`=`$ $`{\displaystyle \frac{g^{}}{1\left(1\frac{g^{}}{g}\right)\left(\frac{t_o}{t}\right)^{\epsilon /2}}}`$ (37) The running coupling $`g(t,t_o)`$ is the effective RG flow of the reaction rate. It can be verified that $`lim_t\mathrm{}g(t,t_o)=g^{}`$, which is of order $`\epsilon `$. Hence for large times we can convert the loop expansion to an $`\epsilon `$-expansion. To obtain answers in $`d=2`$ we take the limit $`\epsilon 0`$ in equations (36), (37): $`M_{N,R}(t,g(t,t_o))`$ $`=`$ $`\left({\displaystyle \frac{t_o}{t}}\right)^Nn^N(t_o,g(t,t_o))`$ (38) $`\left[{\displaystyle \frac{g(t,t_o)}{g}}\right]^{\frac{N(N1)}{2}}`$ $`g(t,t_o)`$ $`=`$ $`{\displaystyle \frac{g}{1+\frac{g}{4\pi }\mathrm{ln}\left(\frac{t}{t_o}\right)}}`$ (39) For large times $`g(t,t_o)\frac{4\pi }{\mathrm{ln}\left(\frac{t}{t_o}\right)}`$. Then in the large time limit we obtain the following scaling in $`t`$ (recall $`M_{N,R}(t)\varphi ^N(t)_R`$) : $$M_{N,R}(t)\{\begin{array}{cc}t^{Nd/2}t^{\frac{N(N1)\epsilon }{4}+𝒪(\epsilon ^2)}\hfill & d<2\hfill \\ \left(\frac{\mathrm{ln}t}{t}\right)^N\left(\mathrm{ln}t\right)^{\frac{N(N1)}{2}}\left(1+𝒪(\frac{1}{\mathrm{ln}t})\right)\hfill & d=2\hfill \end{array}$$ (40) Combining (14) and (40) we obtain the following results for $`P_t(N,\mathrm{\Delta }V)`$ $$\frac{P_t(N,\mathrm{\Delta }V)}{P_t(1,\mathrm{\Delta }V)^N}\{\begin{array}{cc}\left(\frac{\mathrm{\Delta }V^{2/d}}{t}\right)^{\frac{N(N1)\epsilon }{4}+𝒪(\epsilon ^2)}\hfill & d<2\hfill \\ \left(\mathrm{ln}\left[\frac{t}{\mathrm{\Delta }V}\right]\right)^{\frac{N(N1)}{2}}\left(1+𝒪(\frac{1}{\mathrm{ln}t})\right)\hfill & d=2\hfill \end{array}$$ (41) where the $`\mathrm{\Delta }V`$ dependence is restored using dimensional arguments. The physical interpretation of (41) is that in the large-time limit particles are anti-correlated (recall equation (2)): given the same average density the probability of finding $`N`$ reacting particles in $`\mathrm{\Delta }V`$ goes to zero faster than the probability of finding $`N`$ non-reacting particles in $`\mathrm{\Delta }V`$. The origin of anti-correlation can be traced back to the recurrence property of random walks in $`d2`$. Let us compare (41) with exact results for $`N=1,2`$. For $`N=1`$ there is no anomaly and our formula is in agreement with the well-known result found in Lee for $`d2`$. Let us examine the case $`N=2`$ by studying the Langevin SDE : $`_t\varphi =2\lambda \varphi ^2`$ \+ Ito noise term Oleg2 , which is equivalent to the field theory (8). Taking averages of both sides yields the first Hopf equation: $`_t\varphi =2\lambda \varphi ^2`$. For $`d<2`$ we know $`\varphi t^{d/2}`$. Substituting this into the Hopf equation gives $`\varphi ^2t^{(1+d/2)}=t^{(d+\epsilon /2)}`$. In $`d=2`$, $`\varphi \frac{\mathrm{ln}t}{t}`$ and using this in the Hopf equation yields $`\varphi ^2\frac{\mathrm{ln}t}{t^2}`$. These are exact relations and are in agreement with our formula (40) with $`𝒪(\epsilon ^2)=0`$. We conjecture that for $`N=2`$, the $`𝒪(\epsilon ^2)`$ corrections are absent in (40). Finally, note from equation (23) $`lim_\lambda \mathrm{}g=g^{}`$. Thus for $`d2`$, the limit $`\lambda \mathrm{}`$ for finite $`t`$ yields the same asymptotics as the limit $`t\mathrm{}`$ for finite $`\lambda `$. Therefore the large time asymptotics of the model at hand belongs to the universality class of instantaneous annihilation. ### III.3 Comparison of results of $`\epsilon `$-expansion with exact results in d=1 Equation (41) gives us an asymptotically exact result in $`d=2`$ $`(\epsilon =0)`$. In the previous section we have also confirmed that order-$`\epsilon `$ expansion of scaling exponent of $`P(2,\mathrm{\Delta }V)`$ given by the first equation in (41) is exact in all dimensions $`d2`$. As it turns out, (41) yields an exact answer for multi-scaling of probabilities in $`d=1`$ for all values of $`N`$. We are fortunate to have some exact results for multi-point correlation functions in the problem of diffusion-limited annihilation $`A+A\mathrm{}`$ for $`d=1`$, $`\lambda =\mathrm{}`$ mass2 . For more details please refer to mass2 ; mass1 ; ben . Using Mathematica and recurrence relations derived in mass2 we were able to compute exact analytical expressions for $`P_t(N,\mathrm{\Delta }V)`$ for $`N=1,2,3,4`$. Based on the results we find that the $`𝒪(\epsilon ^2)`$ terms are absent in the $`\epsilon `$-expansion (41). This leads us to conjecture that in $`d=1`$ the one-loop answer for the scaling exponents are exact. To substantiate our claims, let us review the key results from mass2 . The reaction rate is infinite, hence we do not expect to find more than one particle at a given site. We will use the notation used in mass2 ; mass1 ; ben . The correlation function $`\rho _N(x_1,\mathrm{},x_N;t)`$ represents the joint probability density of finding N particles positioned at $`x_1,\mathrm{},x_N`$ at time $`t`$. In particular, $`\rho (t)\rho _1(t)`$ is the average density. There is also the convention that $`x_i<x_j`$ for $`i<j`$. In the limit of large times: $$P_t(N,\mathrm{\Delta }V)=_{\mathrm{\Delta }V}𝑑x_1\mathrm{}𝑑x_N\rho _N(x_1,\mathrm{},x_N;t)$$ (42) In the large time limit the following answers hold true: $`\rho (t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{8\pi Dt}}}`$ (43) $`{\displaystyle \frac{\rho _2(x_1,x_2;t)}{\rho ^2(t)}}`$ $`=`$ $`1e^{2z_{21}^2}+\sqrt{\pi }z_{21}e^{z_{21}^2}\mathrm{erfc}(z_{21})`$ (44) $`z_{ji}`$ $`=`$ $`{\displaystyle \frac{x_jx_i}{\sqrt{8Dt}}},`$ (45) where $`D`$ is the diffusion coefficient mass2 . Note that as $`x_2x_1`$ (or vice-versa), the correlation function vanishes. This is a reflection of anti-correlations between particles. For small separations, it can be easily shown that $$\frac{\rho _2(x_1,x_2;t)}{\rho ^2(t)}=\sqrt{\pi }z_{21}+𝒪(z_{21}^2)$$ (46) The above expression for $`\rho _2`$ is valid in the limit of large time and fixed separation $`\mathrm{\Delta }=x_2x_1`$, as in this limit $`z_{21}0`$. Due to (42) the scaling for (46) agrees with our answer for $`N=2`$ in (41): $$\frac{P_t(2,\mathrm{\Delta }V)}{P_t(1,\mathrm{\Delta }V)^2}\frac{\mathrm{\Delta }}{t^{1/2}}$$ (47) With the aid of Mathematica and using similar arguments as above, it can be shown that $`\frac{\rho _3(x_1,x_2,x_3;t)}{\rho ^3(t)}t^{3/2}`$ and $`\frac{\rho _4(x_1,x_2,x_3,x_4;t)}{\rho ^4(t)}t^3`$ , which are in agreement with our conjecture for the cases $`N=3`$ and $`N=4`$ respectively. The formal reason for anomalous scaling for $`N=2`$ is vanishing of the two-particle probability distribution function at $`x_1=x_2`$. This is the most clear indication of anti-correlation between annihilating particles. This same phenomenon is responsible for zeros of multi-particle distribution functions as well. To explore the nature of these zeros starting from the exact recurrence relations of mass2 we were forced to use Mathematica: exact expressions for correlation functions are written as linear combination of a variety of terms involving products of $`\mathrm{exp}(z_{ji}^2),\mathrm{erfc}(z_{ji}),z_{ji}`$. At small separations these expressions simplify due to a number of cancellations which are not obvious. Using Mathematica to break through the tedious computations we find: $`{\displaystyle \frac{\rho _2(x_1,x_2;t)}{\rho ^2(t)}}`$ $`=`$ $`\sqrt{\pi }z_{21}`$ (48) $`{\displaystyle \frac{\rho _3(x_1,x_2,x_3;t)}{\rho ^3(t)}}`$ $`=`$ $`6\sqrt{\pi }z_{21}z_{31}z_{32}`$ (49) $`{\displaystyle \frac{\rho _4(x_1,x_2,x_3,x_4;t)}{\rho ^4(t)}}`$ $`=`$ $`2\pi z_{21}z_{31}z_{32}z_{41}z_{42}z_{43}`$ (50) Note that all distribution functions above vanish as the first power of separation between any pair of particles. Hence we can make a conjecture about spatio-temporal behaviour of distribution functions for arbitrary $`N`$: $$\frac{\rho _N(x_1,\mathrm{},x_N;t)}{\rho ^N(t)}\underset{1i<jN}{}z_{ji}$$ (51) The above expression is a direct generalization of (50) based on permutation symmetry and self-similarity. It states that the spatial dependence of the probability density is given by the Van-der-Monde determinant of particles’ coordinates. Conjecture (51) reproduces the temporal scaling derived in (40) as the right hand side contains $`N(N1)/2`$ factors $`z_{ji}`$ each contributing $`t^{1/2}`$ to the scaling. We will present a rigorous proof of this conjecture in a separate publication. ## IV Reaction $`3A\mathrm{}`$ in critical dimension $`d=1`$ In the final section of this paper, we extend the preceding analysis to the problem $`3A\mathrm{}`$, where reactions occur in triples. Unlike the binary case we cannot apply EIM to analyze this reaction in $`d=1`$. Thus we will adopt the field-theoretic approach. We give a brief introduction to the important quantities in the model Lee . The action $`S`$ given by $`S`$ $`=`$ $`{\displaystyle d^dx_0^\tau 𝑑t\overline{\varphi }(_t\mathrm{\Delta })\varphi }+3\lambda \overline{\varphi }\varphi ^3+3\lambda \overline{\varphi }^2\varphi ^3+`$ (52) $`\lambda \overline{\varphi }^3\varphi ^3n_0\overline{\varphi }\delta (t)`$ The Feynman rules are shown in Fig 4. The action must be dimensionless which requires the following: $$[t]=L^2[\varphi ]=L^d[\overline{\varphi }]=L^0[\lambda ]=L^{2d2}$$ (53) The reaction rate is dimensionless at $`d=1`$, which is the critical dimension for this problem. Hence in $`d=1`$ we expect the system to be characterized by MFT with logarithmic corrections. As for the binary system, we use the classical (MF) versions of the density and response functions to study the loop expansion. The integral equations satisfied by the classical density and classical response function are shown in Fig 5. Their solutions are given by $`n_{cl}(t)`$ $`=`$ $`\left[{\displaystyle \frac{1}{6\lambda t}}\right]^{1/2}`$ (54) $`G_{cl}(x_2,t_2;x_1,t_1)`$ $`=`$ $`\left[{\displaystyle \frac{n_{cl}(t_2)}{n_{cl}(t_1)}}\right]^3G_o(x_2,t_2;x_1,t_1)`$ (55) Power counting shows that reaction rate renormalization takes care of singularities of Greens’ functions which do not contain any composite operators. Moreover it turns out that all divergences of $`\varphi ^2`$ are also eliminated by reaction rate renormalization. For higher order composite operators one needs multiplicative renormalization. Let $`g_o=\lambda t_o^\epsilon `$ be the bare (dimensionless) reaction rate, where $`t_o`$ is our arbitrary reference time. Note that here $`\epsilon =1d`$. To eliminate divergences associated with reaction rate renormalization we need an expression for $`g_o`$ in terms of $`g`$. The relation (23) holds, but $`g^{}=\frac{2\pi \epsilon }{\sqrt{3}}+𝒪(\epsilon ^2)`$. The lowest order diagrams contributing to $`M_N\varphi ^N`$ are shown in Fig 6. The only divergence at the 1-loop level associated with composite operators comes from the connected three-point function. This determines $`Z_N`$ and hence the anomalous scaling. The Callan-Symanzik equation in $`d=1`$ takes the form $$\left[t\frac{}{t}+\beta (g)\frac{}{g}+\frac{N}{2}+\gamma _N(g)\right]M_{N,R}(t,t_o,g)=0$$ (56) where $`\beta (g)`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{2\pi }}g^2`$ (57) $`\gamma _N(g)`$ $`=`$ $`{\displaystyle \frac{N(N1)(N2)}{4\pi \sqrt{3}}}g+𝒪(g^2)`$ (58) Once again we choose the reference time $`t_o`$ small so that the dominant contribution at $`t_o`$ is given by MFT $$M_{N,R}(t_o,g)=n_{cl}^N(t_o,g)$$ (59) The solution of (56) is given by $`M_{N,R}(t,g(t,t_o))`$ $`=`$ $`\left({\displaystyle \frac{t_o}{t}}\right)^{N/2}n_{cl}^N(t_o,g(t,t_o))`$ (60) $`\left[{\displaystyle \frac{g(t,t_o)}{g}}\right]^{\frac{N(N1)(N2)}{6}}`$ $`g(t,t_o)`$ $`=`$ $`{\displaystyle \frac{g}{1+\frac{\sqrt{3}g}{2\pi }\mathrm{ln}\left(\frac{t}{t_o}\right)}}`$ (61) Recall $`M_{N,R}(t)\varphi ^N(t)_R`$. For large times, we obtain the following scaling behaviour of composite operators: $`\varphi ^N(t)_R`$ $``$ $`A_N\left[{\displaystyle \frac{\mathrm{ln}t}{t}}\right]^{N/2}(\mathrm{ln}t)^{\frac{N(N1)(N2)}{6}}`$ (62) $`\left(1+𝒪\left({\displaystyle \frac{1}{\sqrt{\mathrm{ln}t}}}\right)\right)`$ where $$A_N=\left(\frac{1}{\sqrt{6}}\right)^N\left(\frac{\sqrt{3}}{2\pi }\right)^{\frac{3NN(N1)(N2)}{6}}$$ (63) Using the relation (14) between composite operators and probabilities $$\frac{P_t(N,\mathrm{\Delta }V)}{P_t(1,\mathrm{\Delta }V)^N}\left(\mathrm{ln}\left[\frac{t}{\mathrm{\Delta }V}\right]\right)^{\frac{N(N1)(N2)}{6}}\left(1+𝒪\left(\frac{1}{\sqrt{\mathrm{ln}t}}\right)\right)$$ (64) The result (64) reflects the anti-correlation of particles as stated in equation (2). We know from Lee that $`\varphi \left[\frac{\mathrm{ln}t}{t}\right]^{1/2}`$ which agrees with (62) upon setting $`N=1`$. There is no anti-correlation between pairs of particles as reactions only occur in triples. Therefore the anomaly should vanish for $`N=2`$ which agrees with (62),(64). Due to the absence of singularities associated with $`\varphi ^2`$ we know the dominant contribution comes from the disconnected $`n_{cl}`$ diagrams for $`N=2`$. Then $`\varphi ^2\frac{\mathrm{ln}t}{t}`$ Lee which is also in agreement with (62). Finally to check (62) for $`N=3`$, we can use the first Hopf equation for the theory: $`_t\varphi =3\lambda \varphi ^3`$. As $`\varphi \left[\frac{\mathrm{ln}t}{t}\right]^{1/2}`$, the Hopf equation implies that $`\varphi ^3t^{3/2}(\mathrm{ln}t)^{1/2}`$ exactly. This agrees with (62) as well. ## V Summary The main finding of this paper was the multi-scaling of the probability distributions of multi-particle configurations for single species reaction-diffusion systems. The scaling was indicative of particles being anti-correlated in the large time-limit. In particular, the quadratic scaling exponent in the binary system reflects pairwise anti-correlation. For the ternary case, the scaling exponent is cubic which shows anti-correlation within particle triples. We obtained our results in a field-theoretic setting by identifying probability distributions of multi-particle configurations, at scales much smaller than correlation length, with composite operators in Doi-Zeldovich field theory. The origin of the multi-scaling can therefore be traced back to the anomalous dimensions of the corresponding composite operators. We obtained exact logarithmic corrections to scaling for the binary system in $`d=2`$ and for the ternary reaction in $`d=1`$. We computed scaling exponents for the binary system in $`d<2`$ using $`\epsilon `$-expansion. By analyzing the first Hopf equation for the binary system, we proved that the one-loop $`\epsilon `$-expansion gives the exact answer for the probability of finding two particles in a fixed volume. A similar computation for the ternary system confirms the result of the RG computation for the probability of finding three particles in the fixed volume. RG analysis led us to several conjectures for scaling exponents for the $`A+A\mathrm{}`$ system in $`d=1`$. First, by comparing one-loop answers with the exact results in one dimension mass2 for $`N=1,2,3,4`$, we conjecture that two and higher loop corrections are absent in $`d=1`$ for an arbitrary number of particles $`N`$, see (41). Second, basing on the previous conjecture and dimensional analysis we propose that the spatial dependence of the multi-particle probability density is given by the absolute value of the Van-der-Monde determinant of particles’ positions, see (51). We have recently found a rigorous proof of the stated conjectures, which will be published separately. ## VI Acknowledgments We would like to thank Colm Connaughton for useful discussions and help with Mathematica, and Roger Tribe for numerous useful discussions. R. M. gratefully acknowledges MIND for support of this research.
warning/0506/nucl-th0506039.html	ar5iv	text	# Two-photon exchange in elastic electron-nucleon scattering ## I Introduction Electromagnetic form factors are fundamental observables which characterize the composite nature of the nucleon. Several decades of elastic form factor experiments with electron beams, including recent high-precision measurements at Jefferson Lab and elsewhere, have provided considerable insight into the detailed structure of the nucleon. In the standard one-photon exchange (Born) approximation, the electromagnetic current operator is parameterized in terms of two form factors, usually taken to be the Dirac ($`F_1`$) and Pauli ($`F_2`$) form factors, $`\mathrm{\Gamma }^\mu =F_1(q^2)\gamma ^\mu +{\displaystyle \frac{i\sigma ^{\mu \nu }q_\nu }{2M}}F_2(q^2),`$ (1) where $`q`$ is the momentum transfer to the nucleon, and $`M`$ is the nucleon mass. The resulting cross section depends on two kinematic variables, conventionally taken to be $`Q^2q^2`$ (or $`\tau Q^2/4M^2`$) and either the scattering angle $`\theta `$, or the virtual photon polarization $`\epsilon =\left(1+2(1+\tau )\mathrm{tan}^2(\theta /2)\right)^1`$. In terms of the Sachs electric and magnetic form factors, defined as $`G_E(Q^2)`$ $`=`$ $`F_1(Q^2)\tau F_2(Q^2),`$ (2) $`G_M(Q^2)`$ $`=`$ $`F_1(Q^2)+F_2(Q^2),`$ (3) the reduced Born cross section can be written $`\sigma _R=G_M^2(Q^2)+{\displaystyle \frac{\epsilon }{\tau }}G_E^2(Q^2).`$ (4) The standard method which has been used to determine the electric and magnetic form factors, particularly those of the proton, has been the Rosenbluth, or longitudinal-transverse (LT), separation method. Since the form factors in Eq. (4) are functions of $`Q^2`$ only, studying the cross section as a function of the polarization $`\epsilon `$ at fixed $`Q^2`$ allows one to extract $`G_M^2`$ from the $`\epsilon `$-intercept, and the ratio $`R\mu G_E/G_M`$ from the slope in $`\epsilon `$, where $`\mu `$ is the nucleon magnetic moment. The results of the Rosenbluth measurements for the proton have generally been consistent with $`R1`$ for $`Q^26`$ GeV<sup>2</sup> Wal94 ; Arr03 ; Chr04 . The “Super-Rosenbluth” experiment at Jefferson Lab Qat04 , in which smaller systematic errors were achieved by detecting the recoiling proton rather than the electron, as in previous measurements, is also consistent with the earlier LT results. An alternative method of extracting the ratio $`R`$ has been developed recently at Jefferson Lab Jon00 , in which a polarized electron beam scatters from an unpolarized target, with measurement of the polarization of the recoiling proton. From the ratio of the transverse to longitudinal recoil polarizations one finds $`R`$ $`=`$ $`\mu {\displaystyle \frac{E_1+E_3}{2M}}\mathrm{tan}{\displaystyle \frac{\theta }{2}}{\displaystyle \frac{P_T}{P_L}}=\mu \sqrt{{\displaystyle \frac{\tau (1+\epsilon )}{2\epsilon }}}{\displaystyle \frac{P_T}{P_L}},`$ (5) where $`E_1`$ and $`E_3`$ are the initial and final electron energies, and $`P_T`$ ($`P_L`$) is the polarization of the recoil proton transverse (longitudinal) to the proton momentum in the scattering plane. The polarization transfer experiments yielded strikingly different results compared with the LT separation, with $`R10.135(Q^2/\mathrm{GeV}^20.24)`$ over the same range in $`Q^2`$ Arr03 . Recall that in perturbative QCD one expects $`F_1Q^2F_2`$ at large $`Q^2`$ (or equivalently $`G_EG_M`$) pQCD , so that these results imply a strong violation of scaling behavior (see also Refs. Ral03 ; Bel03 ). The question of which experiments are correct has been debated over the past several years. Attempts to reconcile the different measurements have been made by several authors Gui03 ; Blu03 ; Che04 ; Afa05 , who considered whether 2$`\gamma `$ exchange effects, which form part of the radiative corrections (RCs), and which are treated in an approximate manner in the standard RC calculations MT69 , could account for the observed discrepancy. An explicit calculation Blu03 of the two-photon exchange diagram, in which nucleon structure effects were for the first time fully incorporated, indeed showed that around half of the discrepancy could be removed just by the nucleon elastic intermediate states. A partonic level calculation Che04 ; Afa05 subsequently showed that the deep inelastic region can also contribute significantly to the box diagram. In this paper we further develop the methodology introduced in Ref. Blu03 , and apply it to systematically calculate the 2$`\gamma `$ exchange effects in a number of electron–nucleon scattering observables. We focus on the nucleon elastic intermediate states; inelastic contributions are discussed elsewhere Kon05 . In Sec. II we examine the effects of 2$`\gamma `$ exchange on the ratio of electric to magnetic form factors in unpolarized scattering. In contrast to the earlier analysis Blu03 , in which simple monopole form factors were utilized at the internal $`\gamma NN`$ vertices, here we parameterize the vertices by realistic form factors, and study the model-dependence of effects on the ratio $`R`$ due to the choice of form factors. We also compare the results with data on the ratio of $`e^+p`$ to $`e^{}p`$ scattering cross sections, which is directly sensitive to 2$`\gamma `$ exchange effects. In Sec. III we examine the effects of 2$`\gamma `$ exchange on the polarization transfer reaction, $`\stackrel{}{e}pe\stackrel{}{p}`$, for both longitudinally and transversely polarized recoil protons. We also consider the case of proton polarization normal to the reaction plane, which depends on the imaginary part of the box diagram. Since this is absent in the Born approximation, the normal polarization provides a clean signature of 2$`\gamma `$ exchange effects, even though it does not directly address the $`G_E^p/G_M^p`$ discrepancy. Following the discussion of the proton, in Sec. IV we consider 2$`\gamma `$ exchange corrections to the form factors of the neutron, both for the LT separation and polarization transfer techniques. Applying the same formalism to the case of the <sup>3</sup>He nucleus, in Sec. V we compute the elastic contribution from the box diagram to the ratio of charge to magnetic form factors of <sup>3</sup>He. In Sec. VI we summarize our findings, and discuss future work. ## II Two-photon exchange in unpolarized scattering In this section we outline the formalism used to calculate the 2$`\gamma `$ exchange contribution to the unpolarized electron–nucleon cross section, and examine the effect on the $`G_E^p/G_M^p`$ ratio extracted using LT separation. Since there are in general three form factors that are needed to describe elastic $`eN`$ scattering beyond 1$`\gamma `$ exchange, we also evaluate the 2$`\gamma `$ contributions to each of the form factors separately. In the final part of this section, we examine the effect of the 2$`\gamma `$ correction on the ratio of $`e^+p`$ to $`e^{}p`$ elastic cross sections, which is directly sensitive to 2$`\gamma `$ exchange effects. ### II.1 Formalism For the elastic scattering process we define the momenta of the initial electron and nucleon as $`p_1`$ and $`p_2`$, and of the final electron and nucleon as $`p_3`$ and $`p_4`$, respectively, $`e(p_1)+p(p_2)e(p_3)+p(p_4)`$. The four-momentum transferred from the electron to the nucleon is given by $`q=p_4p_2=p_1p_3`$ (with $`Q^2q^2>0`$), and the total electron and proton invariant mass squared is given by $`s=(p_1+p_2)^2=(p_3+p_4)^2`$. In the Born approximation, the amplitude can be written $$_0=i\frac{e^2}{q^2}\overline{u}(p_3)\gamma _\mu u(p_1)\overline{u}(p_4)\mathrm{\Gamma }^\mu (q)u(p_2),$$ (6) where $`e`$ is the electron charge, and $`\mathrm{\Gamma }^\mu `$ is given by Eq. (1). In terms of the amplitude $`_0`$, the corresponding differential Born cross section is given by $`{\displaystyle \frac{d\sigma _0}{d\mathrm{\Omega }}}`$ $`=`$ $`\left({\displaystyle \frac{\alpha }{4Mq^2}}{\displaystyle \frac{E_3}{E_1}}\right)^2\left\|_0\right\|^2=\sigma _{\mathrm{Mott}}{\displaystyle \frac{\tau }{\epsilon (1+\tau )}}\sigma _R,`$ (7) where $`\sigma _R`$ is the reduced cross section given in Eq. (4), and the Mott cross section for the scattering from a point particle is $`\sigma _{\mathrm{Mott}}`$ $`=`$ $`{\displaystyle \frac{\alpha ^2E_3\mathrm{cos}^2\frac{\theta }{2}}{4E_1^3\mathrm{sin}^4\frac{\theta }{2}}},`$ (8) with $`E_1`$ and $`E_3`$ the initial and final electron energies, and $`\alpha =e^2/4\pi `$ the electromagnetic fine structure constant. Including radiative corrections to order $`\alpha `$, the elastic scattering cross section is modified as $`{\displaystyle \frac{d\sigma _0}{d\mathrm{\Omega }}}`$ $``$ $`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}(1+\delta ),`$ (9) where $`\delta `$ includes one-loop virtual corrections (vacuum polarization, electron and proton vertex, and two photon exchange corrections), as well as inelastic bremsstrahlung for real photon emission MT69 . According to the LT separation technique, one extracts the ratio $`R^2`$ from the $`\epsilon `$ dependence of the cross section at fixed $`Q^2`$. Because of the factor $`\epsilon /\tau `$ multiplying $`G_E^2`$ in Eq. (4), the cross section becomes dominated by $`G_M^2`$ with increasing $`Q^2`$, while the relative contribution of the $`G_E^2`$ term is suppressed. Hence understanding the $`\epsilon `$ dependence of the radiative correction $`\delta `$ becomes increasingly important at high $`Q^2`$. As pointed out in Ref. Arr03 , for example, a few percent change in the $`\epsilon `$ slope in $`d\sigma `$ can lead to a sizable effect on $`R`$. In contrast, as we discuss in Sec. III below, the polarization transfer technique does not show the same sensitivity to the $`\epsilon `$ dependence of $`\delta `$. If we denote the amplitude for the one-loop virtual corrections by $`_1`$, then $`_1`$ can be written as the sum of a “factorizable” term, proportional to the Born amplitude $`_0`$, and a non-factorizable part $`\overline{}_1`$, $$_1=f(Q^2,\epsilon )_0+\overline{}_1.$$ (10) The ratio of the full cross section (to order $`\alpha `$) to the Born can therefore be written as $`1+\delta `$ $`=`$ $`{\displaystyle \frac{\left\|_0+_1\right\|^2}{\left\|_0\right\|^2}},`$ (11) with $`\delta `$ given by $$\delta =2f(Q^2,\epsilon )+\frac{2e\{_0^{}\overline{}_1\}}{\|_0\|^2}.$$ (12) In practice the factorizable terms parameterized by $`f(Q^2,\epsilon )`$, which includes the electron vertex correction, vacuum polarization, and the infrared (IR) divergent parts of the nucleon vertex and two-photon exchange corrections, are found to be dominant. Furthermore, these terms are all essentially independent of hadronic structure. However, as explained in Ref. Blu03 , the contributions to the functions $`f(Q^2,\epsilon )`$ from the electron vertex, vacuum polarization, and proton vertex terms depend only on $`Q^2`$, and therefore have no relevance for the LT separation aside from an overall normalization factor. Hence, of the factorizable terms, only the IR divergent two-photon exchange contributes to the $`\epsilon `$ dependence of the virtual photon corrections. The terms which do depend on hadronic structure are contained in $`\overline{}_1`$, and arise from the finite nucleon vertex and two-photon exchange corrections. For the case of the proton, the hadronic vertex correction was analyzed by Maximon and Tjon MT00 , and found to be $`<0.5\%`$ for $`Q^2<6`$ GeV<sup>2</sup>. Since the proton vertex correction does not have a strong $`\epsilon `$ dependence, it will not affect the LT analysis, and can be safely neglected. For the inelastic bremsstrahlung cross section, the amplitude for real photon emission can also be written in the form of Eq. (10). In the soft photon approximation the amplitude is completely factorizable. A significant $`\epsilon `$ dependence arises due to the frame dependence of the angular distribution of the emitted photon. These corrections, together with external bremsstrahlung, contain the main $`\epsilon `$ dependence of the radiative corrections, and are usually accounted for in the experimental analyses. They are generally well understood, and in fact enter differently depending on whether the electron or proton are detected in the final state. Hence corrections beyond the standard $`𝒪(\alpha )`$ radiative corrections which can lead to non-negligible $`\epsilon `$ dependence are confined to the 2$`\gamma `$ exchange diagrams, illustrated in Fig. 1, and are denoted by $`^{2\gamma }`$, which we will focus on in the following. The 2$`\gamma `$ exchange correction $`\delta ^{2\gamma }`$ which we calculate is then essentially $`\delta ^{2\gamma }`$ $``$ $`{\displaystyle \frac{2e\left\{_0^{}^{2\gamma }\right\}}{\left\|_0\right\|^2}}.`$ (13) In principle the two-photon exchange amplitude $`^{2\gamma }`$ includes all possible hadronic intermediate states in Fig. 1. Here we consider only the elastic contribution to the full response function, and assume that the proton propagates as a Dirac particle (excited state contributions are considered in Ref. Kon05 ). We also assume that the structure of the off-shell current operator is similar to that in Eq. (1), and use phenomenological form factors at the $`\gamma NN`$ vertices. This is of course the source of the model dependence in the problem. Clearly this also creates a tautology, as the radiative corrections are also used to determine the experimental form factors. However, because $`\delta `$ is a ratio, the model dependence cancels somewhat, provided the same phenomenological form factors are used for both $`_0`$ and $`^{2\gamma }`$ in Eq. (13). The total 2$`\gamma `$ exchange amplitude, including the box and crossed box diagrams in Fig. 1, has the form $`^{2\gamma }`$ $`=`$ $`e^4{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{N_{\mathrm{box}}(k)}{D_{\mathrm{box}}(k)}}+e^4{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{N_{\mathrm{x}\mathrm{box}}(k)}{D_{\mathrm{x}\mathrm{box}}(k)}},`$ (14) where the numerators are the matrix elements $`N_{\mathrm{box}}(k)`$ $`=`$ $`\overline{u}(p_3)\gamma _\mu (\text{ /}p_1\text{ /}k+m)\gamma _\nu u(p_1)`$ (15) $`\times `$ $`\overline{u}(p_4)\mathrm{\Gamma }^\mu (qk)(\text{ /}p_2+\text{ /}k+M)\mathrm{\Gamma }^\nu (k)u(p_2),`$ $`N_{\mathrm{x}\mathrm{box}}(k)`$ $`=`$ $`\overline{u}(p_3)\gamma _\nu (\text{ /}p_3+\text{ /}k+m)\gamma _\mu u(p_1)`$ (16) $`\times `$ $`\overline{u}(p_4)\mathrm{\Gamma }^\mu (qk)(\text{ /}p_2+\text{ /}k+M)\mathrm{\Gamma }^\nu (k)u(p_2),`$ and the denominators are products of propagators $`D_{\mathrm{box}}(k)`$ $`=`$ $`[k^2\lambda ^2][(kq)^2\lambda ^2]`$ (17) $`\times [(p_1k)^2m^2][(p_2+k)^2M^2],`$ $`D_{\mathrm{x}\mathrm{box}}(k)`$ $`=`$ $`D_{\mathrm{box}}(k)\|_{p_1kp_3+k}.`$ (18) An infinitesimal photon mass $`\lambda `$ has been introduced in the photon propagator to regulate the IR divergences. The IR divergent part is of interest since it is the one usually included in the standard RC analyses. The finite part, which is typically neglected, has been included in Ref. Blu03 and found to have significant $`\epsilon `$ dependence. The IR divergent part of the amplitude $`^{2\gamma }`$ can be separated from the IR finite part by analyzing the structure of the photon propagators in the integrand of Eq. (14). The two poles, where the photons are soft, occur at $`k=0`$ and at $`k=q`$. The dominant (IR divergent) contribution to the integral (14) comes from the poles, and one therefore typically makes the approximation $$_{\mathrm{IR}}^{2\gamma }e^4N_{\mathrm{box}}(0)\frac{d^4k}{(2\pi )^4}\frac{1}{D_{\mathrm{box}}(k)}+e^4N_{\mathrm{x}\mathrm{box}}(0)\frac{d^4k}{(2\pi )^4}\frac{1}{D_{\mathrm{x}\mathrm{box}}(k)},$$ (19) with $`N_{\mathrm{box}}(q)`$ $`=`$ $`N_{\mathrm{box}}(0)=4ip_1p_2{\displaystyle \frac{q^2_0}{e^2}},`$ (20) $`N_{\mathrm{x}\mathrm{box}}(q)`$ $`=`$ $`N_{\mathrm{x}\mathrm{box}}(0)=4ip_3p_2{\displaystyle \frac{q^2_0}{e^2}}.`$ (21) In this case the IR divergent contribution is proportional to the Born amplitude, and the corresponding correction to the Born cross section is independent of hadronic structure. The remaining integrals over propagators can be done analytically. In the target rest frame the total IR divergent two-photon exchange contribution to the cross section is found to be $$\delta _{\mathrm{IR}}=\frac{2\alpha }{\pi }\mathrm{ln}\left(\frac{E_1}{E_3}\right)\mathrm{ln}\left(\frac{Q^2}{\lambda ^2}\right),$$ (22) a result given by Maximon and Tjon MT00 . The logarithmic IR singularity in $`\lambda `$ is exactly cancelled by a corresponding term in the bremsstrahlung cross section involving the interference between real photon emission from the electron and from the nucleon. By contrast, in the standard treatment of Mo and Tsai (MT) MT69 a different approximation for the integrals over propagators is introduced. Here, the IR divergent contribution to the cross section is $$\delta _{\mathrm{IR}}(\mathrm{MT})=2\frac{\alpha }{\pi }\left(K(p_1,p_2)K(p_3,p_2)\right),$$ (23) where $`K(p_i,p_j)=p_ip_j_0^1𝑑y\mathrm{ln}(p_y^2/\lambda ^2)/p_y^2`$ and $`p_y=p_iy+p_j(1y)`$. The logarithmic dependence on $`\lambda `$ is the same as Eq. (22), however. As mentioned above, the full expression in Eq. (14) includes both finite and IR divergent terms, and form factors at the $`\gamma NN`$ vertices. In Ref. Blu03 the proton form factors $`F_1`$ and $`F_2`$ were expressed in terms of the Sachs electric and magnetic form factors, $`F_1(Q^2)`$ $`=`$ $`{\displaystyle \frac{G_E(Q^2)+\tau G_M(Q^2)}{1+\tau }},`$ (24) $`F_2(Q^2)`$ $`=`$ $`{\displaystyle \frac{G_M(Q^2)G_E(Q^2)}{1+\tau }},`$ (25) with $`G_E`$ and $`G_M`$ both parameterized by a simple monopole form, $`G_{E,M}(Q^2)\mathrm{\Lambda }^2/(\mathrm{\Lambda }^2+Q^2)`$, with the mass parameter $`\mathrm{\Lambda }`$ related to the size of the proton. In the present analysis we generalize this approach by using more realistic form factors in the loop integration, consistent with the actual $`G_{E,M}`$ data. The functions $`F_1`$ and $`F_2`$ are parameterized directly in terms of sums of monopoles, of the form $`F_{1,2}(Q^2)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{n_i}{d_i+Q^2}},`$ (26) where $`n_i`$ and $`d_i`$ are free parameters, and $`n_N`$ is determined from the normalization condition, $`n_N=d_N(F_{1,2}(0)_{i=1}^{N1}n_i/d_i)`$. The parameters $`n_i`$ and $`d_i`$ for the $`F_1`$ and $`F_2`$ form factors of the proton and neutron are given in Table I. The normalization conditions are $`F_1^p(0)=1`$ and $`F_2^p(0)=\kappa _p`$ for the proton, and $`F_1^n(0)=0`$ and $`F_2^n(0)=\kappa _n`$ for the neutron, where $`\kappa _p=1.793`$ and $`\kappa _n=1.913`$ are the proton and neutron anomalous magnetic moments, respectively. In practice we use the parameterization from Ref. Mer96 , and fit the parameterized form factors a sum of three monopoles, except for $`F_2^n`$, which is fitted with $`N=2`$. As discussed in the next section, the sensitivity of the results to the choice of form factor is relatively mild. Of course, one should note that the data to which the form factors are fitted were extracted under the assumption of 1$`\gamma `$ exchange, so that in principle one should iterate the data extraction and fitting procedure for self-consistency. However, within the accuracy of the data and of the 2$`\gamma `$ calculation the effect of this will be small. To obtain the radiatively corrected cross section for unpolarized electron scattering the polarizations of the incoming and outgoing electrons and nucleons in Eqs. (15) and (16) need to be averaged and summed, respectively. The resulting expression involves a product of traces in the Dirac spaces of the electron and nucleon. The trace algebra is tedious but straightforward. It was carried out using the algebraic program FORM Vermaseren and verified independently using the program Tracer Tracer . We also used two independent Mathematica packages (FeynCalc feyncalc and FormCalc formcalc ) to carry out the loop integrals. The packages gave distinct but equivalent analytic expressions, which gave identical numerical results. The loop integrals in Eq. (14) can be expressed in terms of four-point Passarino-Veltman functions PV79 , which have been calculated using Spence function HV79 as implemented by Veltman Veltman . In the actual calculations we have used the FF program ff . The results of the proton calculation are presented in the following section. ### II.2 2$`\gamma `$ Corrections to Proton Form Factors In typical experimental analyses of electromagnetic form factor data Wal94 radiative corrections are implemented using the prescription of Ref. MT69 , including using Eq. (23) to approximate the 2$`\gamma `$ contribution. To investigate the effect of our results on the data analyzed in this manner, we will therefore compare the $`\epsilon `$ dependence of the full calculation with that of $`\delta _{\mathrm{IR}}(\mathrm{MT})`$. To make the comparison meaningful, we will consider the difference $`\mathrm{\Delta }`$ $``$ $`\delta _{\mathrm{full}}\delta _{\mathrm{IR}}(\mathrm{MT}),`$ (27) in which the IR divergences cancel, and which is independent of $`\lambda `$. The results for the difference $`\mathrm{\Delta }`$ between the full calculation and the MT approximation are shown in Fig. 2 for several values of $`Q^2`$ from 1 to 6 GeV<sup>2</sup>. The additional corrections are most significant at low $`\epsilon `$, and essentially vanish at large $`\epsilon `$. At the lower $`Q^2`$ values $`\mathrm{\Delta }`$ is approximately linear in $`\epsilon `$, but significant deviations from linearity are observed with increasing $`Q^2`$, especially at smaller $`\epsilon `$. In Fig. 3(a) we illustrate the model dependence of the results by comparing the results in Fig. 2 at $`Q^2=1`$ and 6 GeV<sup>2</sup> with those obtained using a dipole form for the $`F_1^p`$ and $`F_2^p`$ form factors, with mass $`\mathrm{\Lambda }=0.84`$ GeV. At the lower, $`Q^2=1`$ GeV<sup>2</sup>, value the model dependence is very weak, with essentially no change at all in the slope. For the larger value $`Q^2=6`$ GeV<sup>2</sup> the differences are slightly larger, but the general trend of the correction remains unchanged. We can conclude therefore that the model dependence of the calculation is quite modest. Also displayed is the correction at $`Q^2=12`$ GeV<sup>2</sup>, which will be accessible in future experiments, showing significant deviations from linearity over the entire $`\epsilon `$ range. The results are also relatively insensitive to the high-$`Q^2`$ behavior of the $`G_E^p/G_M^p`$ ratio, as Fig. 3(b) illustrates. Here the correction $`\mathrm{\Delta }`$ is shown at $`Q^2=6`$ GeV<sup>2</sup> calculated using various form factor inputs, from parameterizations obtained by fitting only the LT-separated data Mer96 ; Arr04 , and those in which $`G_E^p`$ is constrained by the polarization transfer data Arr04 ; Bra02 . The various curves are almost indistinguishable, and the dependence on the form factor inputs at lower $`Q^2`$ is expected to be even weaker than that in Fig. 3(b). The effect of the 2$`\gamma `$ corrections on the cross sections can be seen in Fig. 4, where the reduced cross section $`\sigma _R`$, scaled by the square of the dipole form factor, $`G_D`$ $`=`$ $`\left(1+{\displaystyle \frac{Q^2}{0.71\mathrm{GeV}^2}}\right)^2,`$ (28) is plotted as a function of $`\epsilon `$ for several fixed values of $`Q^2`$. In Fig. 4(a) the results are compared with the SLAC data And94 at $`Q^2=3.25`$, 4, 5 and 6 GeV<sup>2</sup>, and with data from the “Super-Rosenbluth” experiment at JLab Qat04 in Fig. 4 (b). In both cases the Born level results (dotted curves), which are obtained using the form factor parameterization of Ref. Bra02 in which $`G_E^p`$ is fitted to the polarization transfer data Jon00 , have slopes which are significantly shallower than the data. With the inclusion of the 2$`\gamma `$ contribution (solid curves), there is a clear increase of the slope, with some nonlinearity evident at small $`\epsilon `$. The corrected results are clearly in better agreement with the data, although do not reproduce the entire correction necessary to reconcile the Rosenbluth and polarization transfer measurements. To estimate the influence of these corrections on the electric to magnetic proton form factor ratio, the simplest approach is to examine how the $`\epsilon `$ slope changes with the inclusion of the 2$`\gamma `$ exchange. Of course, such a simplified analysis can only be approximate since the $`\epsilon `$ dependence is only linear over limited regions of $`\epsilon `$, with clear deviations from linearity at low $`\epsilon `$ and high $`Q^2`$. In the actual data analyses one should apply the correction $`\mathrm{\Delta }`$ directly to the data, as in Fig. 4. However, it is still instructive to obtain an estimate of the effect on $`R`$ by taking the slope over several ranges of $`\epsilon `$. Following Ref. Blu03 , this can be done by fitting the correction $`(1+\mathrm{\Delta })`$ to a linear function of $`\epsilon `$, of the form $`a+b\epsilon `$, for each value of $`Q^2`$ at which the ratio $`R`$ is measured. The corrected reduced cross section in Eq. (4) then becomes $$\sigma _RaG_M^2(Q^2)\left[1+\frac{\epsilon }{\mu ^2\tau }\left(R^2\left[1+\epsilon b/a\right]+\mu ^2\tau b/a\right)\right],$$ (29) where $$R^2=\frac{\stackrel{~}{R}^2\mu ^2\tau b/a}{1+\overline{\epsilon }b/a}$$ (30) is the “true” form factor ratio, corrected for 2$`\gamma `$ exchange effects, and $`\stackrel{~}{R}`$ is the “effective” ratio, contaminated by 2$`\gamma `$ exchange. Note that in Eqs. (29) and (30) we have effectively linearized the quadratic term in $`\epsilon `$ by taking the average value of $`\epsilon `$ (i.e., $`\overline{\epsilon }`$) over the $`\epsilon `$ range being fitted. In contrast to Ref. Blu03 , where the approximation $`a1`$ was made and the quadratic term in $`\epsilon `$ neglected, the use of the full expression in Eq. (30) leads to a small decrease in $`R`$ compared with the approximate form. The shift in $`R`$ is shown in Fig. 5, together with the polarization transfer data. We consider two ranges for $`\epsilon `$: a large range $`\epsilon =0.20.9`$, and a more restricted range $`\epsilon =0.50.8`$. The approximation of linear $`\epsilon `$ dependence of $`\mathrm{\Delta }`$ should be better for the latter, even though in practice experiments typically sample values of $`\epsilon `$ near its lower and upper bounds. A proposed experiment at Jefferson Lab Lin04 aims to test the linearity of the $`\epsilon `$ plot through a precision measurement of the unpolarized elastic cross section. The effect of the 2$`\gamma `$ exchange terms on $`R`$ is clearly significant. As observed in Ref. Blu03 , the 2$`\gamma `$ corrections have the proper sign and magnitude to resolve a large part of the discrepancy between the two experimental techniques. In particular, the earlier results Blu03 using simple monopole form factors found a shift similar to that in for the $`\epsilon =0.50.8`$ range in Fig. 5, which resolves around 1/2 of the discrepancy. The nonlinearity at small $`\epsilon `$ makes the effective slope somewhat larger if the $`\epsilon `$ range is taken between 0.2 and 0.9. The magnitude of the effect in this case is sufficient to bring the LT and polarization transfer points almost to agreement, as indicated in Fig. 5. While the 2$`\gamma `$ corrections clearly play a vital role in resolving most of the form factor discrepancy, it is instructive to understand the origin of the effect on $`R`$ with respect to contributions to the individual $`G_E^p`$ and $`G_M^p`$ form factors. In general the amplitude for elastic scattering of an electron from a proton, beyond the Born approximation, can be described by three (complex) form factors, $`\stackrel{~}{F}_1`$, $`\stackrel{~}{F}_2`$ and $`\stackrel{~}{F}_3`$. The generalized amplitude can be written as Gui03 ; Che04 $``$ $`=`$ $`i{\displaystyle \frac{e^2}{q^2}}\overline{u}(p_3)\gamma _\mu u(p_1)\overline{u}(p_4)\left(\stackrel{~}{F}_1\gamma ^\mu +\stackrel{~}{F}_2{\displaystyle \frac{i\sigma ^{\mu \nu }q_\nu }{2M}}+\stackrel{~}{F}_3{\displaystyle \frac{\gamma KP^\mu }{M^2}}\right)u(p_2),`$ (31) where $`K=(p_1+p_3)/2`$ and $`P=(p_2+p_4)/2`$. The functions $`\stackrel{~}{F}_i`$ (both real and imaginary parts) are in general functions of $`Q^2`$ and $`\epsilon `$. In the 1$`\gamma `$ exchange limit the $`\stackrel{~}{F}_{1,2}`$ functions approach the usual (real) Dirac and Pauli form factors, while the new form factor $`\stackrel{~}{F}_3`$ exists only at the 2$`\gamma `$ level and beyond, $`\stackrel{~}{F}_{1,2}(Q^2,\epsilon )`$ $``$ $`F_{1,2}(Q^2),`$ (32) $`\stackrel{~}{F}_3(Q^2,\epsilon )`$ $``$ $`0.`$ (33) Alternatively, the amplitude can be expressed in terms of the generalized (complex) Sachs electric and magnetic form factors, $`\stackrel{~}{G}_E=G_E+\delta G_E`$ and $`\stackrel{~}{G}_M=G_M+\delta G_M`$, in which case the reduced cross section, up to order $`\alpha ^2`$ corrections, can be written Che04 $`\stackrel{~}{\sigma }_R`$ $`=`$ $`G_M^2+{\displaystyle \frac{\epsilon }{\tau }}G_E^2+2G_M^2e\left\{{\displaystyle \frac{\delta G_M}{G_M}}+\epsilon Y_{2\gamma }\right\}+{\displaystyle \frac{2\epsilon }{\tau }}G_E^2e\left\{{\displaystyle \frac{\delta G_E}{G_E}}+{\displaystyle \frac{G_M}{G_E}}Y_{2\gamma }\right\},`$ (34) where the form factor $`\stackrel{~}{F}_3`$ has been expressed in terms of the ratio $`Y_{2\gamma }`$ $`=`$ $`\stackrel{~}{\nu }{\displaystyle \frac{\stackrel{~}{F}_3}{G_M}},`$ (35) with $`\stackrel{~}{\nu }KP/M^2=\sqrt{\tau (1+\tau )(1+\epsilon )/(1\epsilon )}`$. We should emphasize that the generalized form factors are not observables, and therefore have no intrinsic physical meaning. Thus the magnitude and $`\epsilon `$ dependence of the generalized form factors will depend on the choice of parametrization of the generalized amplitude. For example, the axial parametrization introduces an effective axial vector coupling beyond Born level, and is written as Rek $``$ $`=`$ $`i{\displaystyle \frac{e^2}{q^2}}\{\overline{u}(p_3)\gamma _\mu u(p_1)\overline{u}(p_4)(F_1^{}\gamma ^\mu +F_2^{}{\displaystyle \frac{i\sigma ^{\mu \nu }q_\nu }{2M}})u(p_2)`$ (36) $`+G_A^{}\overline{u}(p_3)\gamma _\mu \gamma _5u(p_1)\overline{u}(p_4)\gamma ^\mu \gamma _5u(p_2)\}.`$ Following Ref. Afa05 , one finds the relationships $`F_1^{}`$ $`=`$ $`\stackrel{~}{F}_1+\stackrel{~}{\nu }\stackrel{~}{F}_3,`$ (37) $`F_2^{}`$ $`=`$ $`\stackrel{~}{F}_2,`$ (38) $`G_A^{}`$ $`=`$ $`\tau \stackrel{~}{F}_3.`$ (39) In Fig. 6 we show the contributions of 2$`\gamma `$ exchange to the (real parts of the) proton $`\stackrel{~}{G}_E`$ and $`\stackrel{~}{G}_M`$ form factors, and the ratio $`Y_{2\gamma }`$ evaluated at $`Q^2=1`$, 3 and 6 GeV<sup>2</sup>. One observes that the 2$`\gamma `$ correction to $`\stackrel{~}{G}_M`$ is large, with a positive slope in $`\epsilon `$ which increases with $`Q^2`$. The correction to $`\stackrel{~}{G}_E`$ is similar to that for $`\stackrel{~}{G}_M`$ at $`Q^2=1`$ GeV<sup>2</sup>, but becomes shallower at intermediate $`\epsilon `$ values for larger $`Q^2`$. Both of these corrections are significantly larger than the $`Y_{2\gamma }`$ correction, which is weakly $`Q^2`$ dependent, and has a small negative slope in $`\epsilon `$ at larger $`Q^2`$. The contribution to $`Y_{2\gamma }`$ is found to be about 5 times smaller than that extracted in phenomenological analyses Gui03 under the assumption that the entire form factor discrepancy is due to the new $`\stackrel{~}{F}_3`$ contribution (see also Ref. Arr05 ). ### II.3 Comparison of $`e^+p`$ to $`e^{}p`$ cross sections Direct experimental evidence for the contribution of 2$`\gamma `$ exchange can be obtained by comparing $`e^+p`$ and $`e^{}p`$ cross sections through the ratio $`R^{e^+e^{}}`$ $``$ $`{\displaystyle \frac{d\sigma ^{(e^+)}}{d\sigma ^{(e^{})}}}`$ (40) $``$ $`{\displaystyle \frac{\left\|_0^{(e^+)}\right\|^2+2e\left\{_0^{(e^+)}^{2\gamma (e^+)}\right\}}{\left\|_0^{(e^{})}\right\|^2+2e\left\{_0^{(e^{})}^{2\gamma (e^{})}\right\}}}.`$ Whereas the Born amplitude $`_0`$ changes sign under the interchange $`e^{}e^+`$, the 2$`\gamma `$ exchange amplitude $`^{2\gamma }`$ does not. The interference of the $`_0`$ and $`^{2\gamma }`$ amplitudes therefore has the opposite sign for electron and positron scattering. Since the finite part of the 2$`\gamma `$ contribution is negative over most of the range of $`\epsilon `$, one would expect to see an enhancement of the ratio of $`e^+`$ to $`e^{}`$ cross sections, $`R^{e^+e^{}}`$ $``$ $`12\mathrm{\Delta },`$ (41) where $`\mathrm{\Delta }`$ is defined in Eq. (27). Although the current data on elastic $`e^{}p`$ and $`e^+p`$ scattering are sparse, there are some experimental constraints from old data taken at SLAC Bro65 ; Mar68 , Cornell And66 , DESY Bar67 and Orsay Bou68 (see also Ref. ArrEE ). The data are predominantly at low $`Q^2`$ and at forward scattering angles, corresponding to large $`\epsilon `$ ($`\epsilon 0.7`$), where the 2$`\gamma `$ exchange contribution is small ($`1\%`$). Nevertheless, the overall trend in the data reveals a small enhancement in $`R^{e^+e^{}}`$ at the lower $`\epsilon `$ values, as illustrated in Fig. 7 (which shows a subset of the data, from the SLAC experiments Bro65 ; Mar68 ). The data in Fig. 7 are compared with our theoretical results, calculated for several fixed values of $`Q^2`$ ($`Q^2=1`$, 3 and 6 GeV<sup>2</sup>). The results are in good agreement with the data, although the errors on the data points are quite large. Clearly better quality data at backward angles, where an enhancement of up to $`10\%`$ is predicted, would be needed for a more definitive test of the 2$`\gamma `$ exchange mechanism. An experiment Bro04 using a beam of $`e^+e^{}`$ pairs produced from a secondary photon beam at Jefferson Lab will make simultaneous measurements of $`e^{}p`$ and $`e^+p`$ elastic cross sections up to $`Q^22`$ GeV<sup>2</sup>. A proposal to perform a precise ($`1\%`$) comparison of $`e^{}p`$ and $`e^+p`$ scattering at $`Q^2=1.6`$ GeV<sup>2</sup> and $`\epsilon 0.4`$ has also been made at the VEPP-3 storage ring VEPP . ## III Polarized electron–proton scattering The results of the 2$`\gamma `$ exchange calculation in the previous section give a clear indication of a sizable correction to the LT-separated data at moderate and large $`Q^2`$. The obvious question which arises is whether, and to what extent, the 2$`\gamma `$ exchange affects the polarization transfer results, which show the dramatic fall-off of the $`G_E^p/G_M^p`$ ratio at large $`Q^2`$. In this section we examine this problem in detail. The polarization transfer experiment involves the scattering of longitudinally polarized electrons from an unpolarized proton target, with the detection of the polarization of the recoil proton, $`\stackrel{}{e}+pe+\stackrel{}{p}`$. (The analogous process whereby a polarized electron scatters elastically from a polarized proton leaving an unpolarized final state gives rise to essentially the same information.) In the Born approximation the spin dependent amplitude is given by $`_0(s_1,s_4)`$ $`=`$ $`i{\displaystyle \frac{e^2}{q^2}}\overline{u}(p_3)\gamma _\mu u(p_1,s_1)\overline{u}(p_4,s_4)\mathrm{\Gamma }^\mu (q)u(p_2),`$ (42) where $`s_1=(s_1^0;\stackrel{}{s}_1)`$ and $`s_4=(s_4^0;\stackrel{}{s}_4)`$ are the spin four-vectors of the initial electron and final proton, respectively, and the spinor $`u(p_1,s_1)`$ is defined such that $`u(p_1,s_1)\overline{u}(p_1,s_1)=(\text{ /}p_1+m)(1+\gamma _5\text{ /}s_1)/2`$, and similarly for $`\overline{u}(p_4,s_4)`$. The spin four-vector (for either the electron or recoil proton) can be written in terms of the 3-dimensional spin vector $`\zeta `$ specifying the spin direction in the rest frame (see e.g. Ref. MP00 ), $`s^\mu `$ $`=`$ $`({\displaystyle \frac{\stackrel{}{\zeta }\stackrel{}{p}}{m}};\stackrel{}{\zeta }+\stackrel{}{p}{\displaystyle \frac{\stackrel{}{\zeta }\stackrel{}{p}}{m(m+E)}}),`$ (43) where $`m`$ and $`E`$ are the mass and energy of the electron or proton. Clearly in the limit $`\stackrel{}{p}0`$, the spin four-vector $`s(0;\stackrel{}{\zeta })`$. Since $`\zeta `$ is a unit vector, one has $`\stackrel{}{\zeta }^2=1`$, and one can verify from Eq. (43) that $`s^2=1`$ and $`ps=0`$. If the incident electron energy $`E_1`$ is much larger than the electron mass $`m`$, the electron spin four-vector can be related to the electron helicity $`h=\stackrel{}{\zeta }_1\stackrel{}{p}_1`$ by $`s_1`$ $``$ $`h{\displaystyle \frac{p_1}{m}}.`$ (44) The coordinate axes are chosen so that the recoil proton momentum $`\stackrel{}{p}_4`$ defines the $`z`$ axis, in which case for longitudinally polarized protons one has $`\stackrel{}{\zeta }=\widehat{p}_4`$. In the 1$`\gamma `$ exchange approximation the elastic cross section for scattering a longitudinally polarized electron with a recoil proton polarized longitudinally is then given by $`{\displaystyle \frac{d\sigma ^{(L)}}{d\mathrm{\Omega }}}`$ $`=`$ $`h\sigma _{\mathrm{Mott}}{\displaystyle \frac{E_1+E_3}{M}}\sqrt{{\displaystyle \frac{\tau }{1+\tau }}}\mathrm{tan}^2{\displaystyle \frac{\theta }{2}}G_M^2.`$ (45) For transverse recoil proton polarization we define the $`x`$ axis to be in the scattering plane, $`\widehat{x}=\widehat{y}\times \widehat{z}`$, where $`\widehat{y}=\widehat{p}_1\times \widehat{p}_3`$ defines the direction perpendicular, or normal, to the scattering plane. The elastic cross section for producing a transversely polarized proton in the final state, with $`\stackrel{}{\zeta }\stackrel{}{p}_4=0`$, is given by $`{\displaystyle \frac{d\sigma ^{(T)}}{d\mathrm{\Omega }}}`$ $`=`$ $`h\sigma _{\mathrm{Mott}}2\sqrt{{\displaystyle \frac{\tau }{1+\tau }}}\mathrm{tan}{\displaystyle \frac{\theta }{2}}G_EG_M.`$ (46) Taking the ratio of the transverse to longitudinal proton cross sections then gives the ratio of the electric to magnetic proton form factors, as in Eq. (5). Note that in the 1$`\gamma `$ exchange approximation the normal polarization is identically zero. The amplitude for the 2$`\gamma `$ exchange diagrams in Fig. 1 with the initial electron and final proton polarized can be written as $`^{2\gamma }(s_1,s_4)`$ $`=`$ $`e^4{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{N_{\mathrm{box}}(k,s_1,s_4)}{D_{\mathrm{box}}(k)}}+e^4{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{N_{\mathrm{x}\mathrm{box}}(k,s_1,s_4)}{D_{\mathrm{x}\mathrm{box}}(k)}},`$ (47) where the numerators are the matrix elements $`N_{\mathrm{box}}(k,s_1,s_4)`$ $`=`$ $`\overline{u}(p_3)\gamma _\mu (\text{ /}p_1\text{ /}k+m)\gamma _\nu u(p_1,s_1)`$ (48) $`\times `$ $`\overline{u}(p_4,s_4)\mathrm{\Gamma }^\mu (qk)(\text{ /}p_2+\text{ /}k+M)\mathrm{\Gamma }^\nu (k)u(p_2),`$ $`N_{\mathrm{x}\mathrm{box}}(k,s_1,s_4)`$ $`=`$ $`\overline{u}(p_3)\gamma _\nu (\text{ /}p_3+\text{ /}k+m)\gamma _\mu u(p_1,s_1)`$ (49) $`\times `$ $`\overline{u}(p_4,s_4)\mathrm{\Gamma }^\mu (qk)(\text{ /}p_2+\text{ /}k+M)\mathrm{\Gamma }^\nu (k)u(p_2),`$ and the denominators are given in Eqs. (17) and (18). The traces in Eqs. (48) and (49) can be evaluated using the explicit expression for the spin-vectors $`s_1`$ and $`s_4`$ in Eqs. (43) and (44). In analogy with the unpolarized case (see Eq. (27)), the spin-dependent corrections to the longitudinal ($`\mathrm{\Delta }_L`$) and transverse ($`\mathrm{\Delta }_T`$) cross sections are defined as the finite parts of the 2$`\gamma `$ contributions relative to the IR expression from Mo & Tsai MT69 in Eq. (23), which are independent of polarization, $`\mathrm{\Delta }_{L,T}`$ $`=`$ $`\delta _{L,T}^{\mathrm{full}}\delta _{\mathrm{IR}}.`$ (50) Experimentally, one does not usually measure the longitudinal or transverse cross section per se, but rather the ratio of the transverse or longitudinal cross section to the unpolarized cross section, denoted $`P_L`$ or $`P_T`$, respectively. Thus the 2$`\gamma `$ exchange correction to the polarization transfer ratio can be incorporated as $$\frac{P_{L,T}^{1\gamma +2\gamma }}{P_{L,T}^{1\gamma }}=\frac{1+\mathrm{\Delta }_{L,T}}{1+\mathrm{\Delta }},$$ (51) where $`\mathrm{\Delta }`$ is the correction to the unpolarized cross section considered in the previous section. The 2$`\gamma `$ exchange contribution relative to the Born term is shown in Fig. 8. The correction to the longitudinal polarization transfer ratio $`P_L`$ is small overall. This is because the correction $`\mathrm{\Delta }_L`$ to the longitudinal cross section is roughly the same as the correction $`\mathrm{\Delta }`$ to the unpolarized cross section. The corrections $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }_L`$ must be exactly the same at $`\theta =180^{}`$ ($`\epsilon =0`$), and our numerical results bear this out. By contrast, the correction to the transverse polarization transfer ratio $`P_T`$ is enhanced at backward angles, and grows with $`Q^2`$. This is due to a combined effect of $`\mathrm{\Delta }_T`$ becoming more positive with increasing $`Q^2`$, and $`\mathrm{\Delta }`$ becoming more negative. In the standard radiative corrections using the results of Mo & Tsai MT69 , the corrections for transverse polarization are the same as those for longitudinal polarization, so that no additional corrections beyond hard bremsstrahlung need be applied MP00 . Because the polarization transfer experiments Jon00 typically have $`\epsilon 0.7`$–0.8, the shift in the polarization transfer ratio in Eq. (5) due to the 2$`\gamma `$ exchange corrections is not expected to be dramatic. If $`R`$ is the corrected (“true”) electric to magnetic form factor ratio, as in Eq. (29), then the measured polarization transfer ratio is $`\stackrel{~}{R}`$ $`=`$ $`R\left({\displaystyle \frac{1+\mathrm{\Delta }_T}{1+\mathrm{\Delta }_L}}\right).`$ (52) Inverting Eq. (52), the shift in the ratio $`R`$ is illustrated in Fig. 9 by the filled circles (offset slightly for clarity). The unshifted results are indicated by the open circles, and the LT separated results are labeled by diamonds. The effect of the 2$`\gamma `$ exchange on the form factor ratio is a very small, $`3\%`$ suppression of the ratio at the larger $`Q^2`$ values, which is well within the experimental uncertainties. Note that the shift in $`R`$ in Eq. (52) does not include corrections due to hard photon bremsstrahlung (which are part of the standard radiative corrections). Since these would make both the numerator and denominator in Eq. (52) even larger, the correction shown in Fig. 9 would represent an upper limit on the shift in $`R`$. Finally, the 2$`\gamma `$ exchange process can give rise to a non-zero contribution to the elastic cross section for a recoil proton polarized normal to the scattering plane. This contribution is purely imaginary, and does not exist in the 1$`\gamma `$ exchange approximation. It is illustrated in Fig. 10, where the ratio $`\mathrm{\Delta }_N`$ of the 2$`\gamma `$ exchange contribution relative to the unpolarized Born contribution is shown as a function of $`\epsilon `$ for several values of $`Q^2`$. (For consistency in notation we denote this correction $`\mathrm{\Delta }_N`$ rather than $`\delta _N`$, even though there is no IR contribution to the normal polarization.) The normal polarization contribution is very small numerically, $`\mathrm{\Delta }_N1\%`$, and has a very weak $`\epsilon `$ dependence. In contrast to $`\mathrm{\Delta }_L`$ and $`\mathrm{\Delta }_T`$, the normal polarization ratio is smallest at low $`\epsilon `$, becoming larger with increasing $`\epsilon `$. Although not directly relevant to the elastic form factor extraction, the observation of protons with normal polarization would provide direct evidence of 2$`\gamma `$ exchange in elastic scattering. Figure 11 shows the normal polarization asymmetry $`A_y`$ as a function of the center of mass scattering angle, $`\mathrm{\Theta }_{\mathrm{cm}}`$, for several values of $`Q^2`$. The asymmetry is relatively small, of the order of 1% at small $`\mathrm{\Theta }_{\mathrm{cm}}`$ for $`Q^23`$ GeV<sup>2</sup>, but grows with $`Q^2`$. The imaginary part of the 2$`\gamma `$ amplitude can also be accessed by measuring the electron beam asymmetry for electrons polarized normal to the scattering plane Wells . Knowledge of the imaginary part of the 2$`\gamma `$ exchange amplitude could be used to constrain models of Compton scattering, although relating this to the real part (as needed for form factor studies) would require a dispersion relation analysis. ## IV Electron–neutron scattering In this section we examine the effect of the 2$`\gamma `$ exchange contribution on the form factors of the neutron. Since the magnitude of the electric form factor of the neutron is relatively small compared with that of the proton, and as we saw in Sec. III the effects on the proton are significant at large $`Q^2`$, it is important to investigate the extent to which $`G_E^n`$ may be contaminated by 2$`\gamma `$ exchange. Using the same formalism as in Secs. II and III, the calculated 2$`\gamma `$ exchange correction for the neutron is shown in Fig. 12 for $`Q^2=1`$, 3 and 6 GeV<sup>2</sup>. Since there is no IR divergent contribution to $`\delta `$ for the neutron, the total 2$`\gamma `$ correction $`\delta ^{\mathrm{full}}`$ is displayed in Fig. 12. In the numerical calculation, the input neutron form factors from Ref. Mer96 are parameterized using the pole fit in Eq. (26), with the parameters given in Table 1. For comparison, the correction at $`Q^2=6`$ GeV<sup>2</sup> is also computed using a 3-pole fit to the form factor parameterization from Ref. Bos95 . The difference between these is an indication of the model dependence of the calculation. The most notable difference with respect to the proton results is the sign and slope of the 2$`\gamma `$ exchange correction. Namely, the magnitude of the correction $`\delta ^{\mathrm{full}}(\epsilon ,Q^2)`$ for the neutron is $`3`$ times smaller than for the proton. The reason for the sign change is the negative anomalous magnetic moment of the neutron. The $`\epsilon `$ dependence is approximately linear at moderate and high $`\epsilon `$, but at low $`\epsilon `$ there exists a clear deviation from linearity, especially at large $`Q^2`$. Translating the $`\epsilon `$ dependence to the form factor ratio, the resulting shift in $`\mu _nG_E^n/G_M^n`$ is shown in Fig. 13 at several values of $`Q^2`$, assuming a linear 2$`\gamma `$ correction over two different $`\epsilon `$ ranges ($`\epsilon =0.20.9`$ and $`\epsilon =0.50.8`$). The baseline (uncorrected) data are from the global fit in Ref. Mer96 . The shift due to 2$`\gamma `$ exchange is small at $`Q^2=1`$ GeV<sup>2</sup>, but increases significantly by $`Q^2=6`$ GeV<sup>2</sup>, where it produces a 50–60% rise in the uncorrected ratio. These results suggest that, as for the proton, the LT separation method is subject to large corrections from 2$`\gamma `$ exchange at large $`Q^2`$. While the 2$`\gamma `$ corrections to the form factor ratio from LT separation are signficant, particularly at large $`Q^2`$, in practice the neutron $`G_E^n`$ form factor is commonly extracted using the polarization transfer method. To compare the 2$`\gamma `$ effects on the ratio $`\mu _nG_E^n/G_M^n`$ extracted by polarization transfer, in Fig. 14 we plot the same “data points” as in Fig. 13, shifted by the $`\delta _{L,T}`$ corrections as in Eq. (52) at two values of $`\epsilon `$ ($`\epsilon =0.3`$ and 0.8). The shift is considerably smaller than that from the LT method, but nevertheless represents an approximately 4% (3%) suppression at $`\epsilon =0.3`$ (0.8) for $`Q^2=3`$ GeV<sup>2</sup>, and $`10\%`$ (5%) suppression for $`Q^2=6`$ GeV<sup>2</sup> for the same $`\epsilon `$. In the Jefferson Lab experiment Madey to measure $`G_E^n/G_M^n`$ at $`Q^2=1.45`$ GeV<sup>2</sup> the value of $`\epsilon `$ was around 0.9, at which the 2$`\gamma `$ correction was $`2.5\%`$. In the recently approved extension of this measurement to $`Q^24.3`$ GeV<sup>2</sup> MadeyNew , the 2$`\gamma `$ correction for $`\epsilon 0.82`$ is expected to be around 3%. While small, these corrections will be important to take into account in order to achieve precision at the several percent level. ## V <sup>3</sup>He Elastic Form Factors In this section we extend our formalism to the case of elastic scattering from <sup>3</sup>He nuclei. Of course, the contribution of <sup>3</sup>He intermediate states in 2$`\gamma `$ exchange is likely to constitute only a part of the entire effect – contributions from break-up channels may also be important. However, we can obtain an estimate on the size of the effect on the <sup>3</sup>He form factors, in comparison with the effect on the nucleon form factor ratio. The expressions used to evaluate the 2$`\gamma `$ contributions are similar to those for the nucleon, since <sup>3</sup>He is a spin-$`\frac{1}{2}`$ particle, although there are some important differences. For instance, the charge is now $`Ze`$ (where $`Z=2`$ is the atomic number of <sup>3</sup>He), the mass $`M_{{}_{}{}^{3}\mathrm{He}}`$ is $`3`$ times larger than the nucleon mass, and the anomalous magnetic moment is $`\kappa _{{}_{}{}^{3}\mathrm{He}}=4.185`$. In addition, the internal $`\gamma {}_{}{}^{3}\mathrm{He}`$ form factor is somewhat softer than the corresponding nucleon form factor (since the charge radius of the <sup>3</sup>He nucleus is $`1.88`$ fm). Using a dipole shape for the form factor gives a cut-off mass of $`\mathrm{\Lambda }_{{}_{}{}^{3}\mathrm{He}}0.37`$ GeV. The 2$`\gamma `$ exchange correction is shown in Fig. 15 as a function of $`\epsilon `$ for several values of $`Q^2`$. The $`\epsilon `$ dependence illustrates the interesting interplay between the Dirac and Pauli contributions to the cross section. At low $`Q^2`$ ($`Q^21`$ GeV<sup>2</sup>), the $`F_1`$ contribution is dominant, and the effect has the same sign and similar magnitude as in the proton. The result in fact reflects a partial cancellation of 2 opposing effects: the larger charge squared $`Z^2`$ of the <sup>3</sup>He nucleus makes the effect larger (by a factor $`4`$), while the larger mass squared of the <sup>3</sup>He nucleus suppresses the effect by a factor $`9`$. In addition, the form factor used is much softer than that of the nucleon, so that the overall effect turns out to be similar in magnitude as for the proton. With increasing $`Q^2`$ the Pauli $`F_2`$ term becomes more important, so that for $`Q^23`$ GeV<sup>2</sup> the overall sign of the contribution is positive. Interestingly, over most of the region between $`\epsilon 0.2`$ and 0.9 the slope in $`\epsilon `$ is approximately constant. This allows us to extract the correction to the ratio of charge to magnetic form factors, $`F_C/F_M`$, which we illustrate in Fig. 16. The effect is a small, $`0.5\%`$ reduction in the ratio for $`Q^23`$ GeV<sup>2</sup>, which turns into an enhancement at large $`Q^2`$. However, the magnitude of the effect is small, and even for $`Q^2=6`$GeV<sup>2</sup> the 2$`\gamma `$ effect only gives $`2\%`$ increase in the form factor ratio. Proposed experiments at Jefferson Lab Pet03 would measure the <sup>3</sup>He form factors to $`Q^24`$ GeV<sup>2</sup>. ## VI Conclusion We have presented a comprehensive analysis of the effects of 2$`\gamma `$ exchange in elastic electron–nucleon scattering, taking particular account of the effects of nucleon structure. Our main purpose has been to quantify the 2$`\gamma `$ effect on the ratio of electric to magnetic form factors of the proton, which has generated controversy recently stemming from conflicting results of measurements at large $`Q^2`$. Consistent with the earlier preliminary investigation Blu03 , we find that inclusion of 2$`\gamma `$ exchange reduces the $`G_E^p/G_M^p`$ ratio extracted from LT-separated cross section data, and resolves a significant amount of the discrepancy with the polarization transfer results. At higher $`Q^2`$ we find strong deviations from linearity, especially at small $`\epsilon `$, which can be tested in future high-precision cross section measurements. There is some residual model-dependence in the calculation of the 2$`\gamma `$ amplitude arising from the choice of form factors at the internal $`\gamma ^{}NN`$ vertices in the loop integration. This dependence, while not overwhelming, will place limitations on the reliability of the LT separation technique in extracting high-$`Q^2`$ form factors. On the other hand, the size of the 2$`\gamma `$ contributions to elastic scattering could be determined from measurement of the ratio of $`e^{}p`$ to $`e^+p`$ elastic cross sections, which are uniquely sensitive to 2$`\gamma `$ exchange effects. We have also generalized our analysis to the case where the initial electron and recoil proton are polarized, as in the polarization transfer experiments. While the 2$`\gamma `$ corrections can be as large as $`4`$–5% at small $`\epsilon `$ for $`Q^26`$ GeV<sup>2</sup>, because the polarization transfer measurements are performed typically at large $`\epsilon `$ we find the impact on the extracted $`G_E^p/G_M^p`$ ratio to be quite small, amounting to $`3\%`$ suppression at the highest $`Q^2`$ value. Extending the formalism to the case of the neutron, we have calculated the 2$`\gamma `$ exchange corrections to the neutron $`G_E^n/G_M^n`$ ratio. While numerically smaller than for the proton, the corrections are nonetheless important since the magnitude of $`G_E^n`$ itself is small compared with $`G_E^p`$. Furthermore, because of the opposite sign of the neutron magnetic moment relative to the proton, the 2$`\gamma `$ corrections to the LT-separated cross section give rise to a sizable enhancement of $`G_E^n/G_M^n`$ at large $`Q^2`$. The analogous effects for the polarization transfer ratio are small, on the other hand, giving rise to a few percent suppression for $`Q^26`$ GeV<sup>2</sup>. Finally, we have also obtained an estimate of the 2$`\gamma `$ exchange contribution to the elastic form factors of <sup>3</sup>He from elastic intermediate states. The results reveal an interesting interplay between an enhancement from the larger charge of the <sup>3</sup>He nucleus and a suppression due to the larger mass. Together with softer form factor (larger radius) compared with that of the nucleon, the net effect is $`1\%`$ over the $`Q^2`$ range accessible to current and upcoming experiments. Contributions from excited states, such as the $`\mathrm{\Delta }`$ and heavier baryons, may modify the quantitative analysis presented here. Naively, one could expect their effect to be suppressed because of the larger masses involved, at least for the real parts of the form factors. An investigation of the inelastic excitation effects is presented in Ref. Kon05 . ###### Acknowledgements. We would like to thank J. Arrington for helpful discussions and communications. This work was supported in part by NSERC (Canada), DOE grant DE-FG02-93ER-40762, and DOE contract DE-AC05-84ER-40150 under which the Southeastern Universities Research Association (SURA) operates the Thomas Jefferson National Accelerator Facility (Jefferson Lab).
warning/0506/math0506450.html	ar5iv	text	# Quantization by cochain twists and nonassociative differentials ## 1. Introduction This paper is a sequel to \[BM1\] in which we studied algebras that were associative to the required order in a deformation parameter $`\mathrm{}`$ but allowed the possibility that the exterior algebra in noncommutative geometry could be nonassociative to that order. We showed that this was necessary for the standard quantum groups $`_q[G]`$, i.e. these associative algebras admit no associative exterior algebra of classical dimensions that is bicovariant. Nonassociative calculi were, however, possible by use of Drinfeld’s twisiting \[D2\] applied in the category of (super)coquasiHopf algebras. In the present work we provide many more examples using not the quasi-Hopf algebra theory itself but a ‘module algebra’ twist theory in which any algebra in the category of modules covariant under the a classical (or quantum) group is also twisted. Such methods have been used to obtain nonassociative algebras\[AM1\] as well as associative ones\[DGM\]. That one obtains differential calculi as well on such algebras is explored in general terms in \[AM2\]. We show now that this setting also allows to obtain associative algebras and induced differential calculi for some very standard and not-quantum-group-related quantizations, but with a similar price to pay. Thus, we use Hopf algebra methods but apply them to classical situations, notably to coadjoint spaces $`𝔤^{}`$ and their quantisation by the enveloping algebra $`U_{\mathrm{}}(𝔤)`$. Clearly this and other ‘noncommutative coordinate algebras’ that we consider are perfectly associative so it is some surprise, and the main result of the present paper, that their natural induced noncommutative differential calculus is again nonassociative. In the case of $`U_{\mathrm{}}(𝔤)`$ we show (Theorem 5.1.2) that any calculus which is translation and $`𝔤`$-covariant and has classical dimensions must be nonassociative. This is analogous to the result in \[BM1\] for quantum groups but now for classical enveloping algebras (and the proof is similar). In this way we confirm and provide major new examples of the general analysis in \[BM1\]. We particularly analyse the semiclassical level of these results in terms of Poisson and symplectic geometry followed by the next-to-semiclassical order. An outline of the paper is as follows. In Section 2 we describe the general algebraic twisting theory that we shall use. Section 3 then describes the special case that will used for all our examples, namely a method of quantisation induced by a classical symmetry and a cochain. Thus we begin with a classical manifold $`M`$ with a classical Lie algebra symmetry group $`\mathrm{diff}(M)`$. As Hopf algebra we take $`H=U()`$ the enveloping algebra. Then the scheme is that any suitable element $`FHH`$ (a cochain) induces a quantisation of $`M`$. We semiclassicalise this theory and see how Poisson-compatible (pre)connections in the sense of \[BM1\] arise out of the choice of $`F`$ and $``$. The choice of the latter covariance Lie algebra determines what kind of connections or preconnections can arise by the cochain twisting construction and hence what structures the quantisation respects. We also briefly discuss the inverse problem of obtaining a cochain $`F`$ and hence a quantisation given a symplectic form and symplectic connection on $`M`$. Section 3.2 analyses the situation for $`M=^{2n}`$ with its standard symplectic structure and general symplectic connection. Sections 4,5,6,7 then turn to the main examples of the paper. These examples are all constructed by a second order or in some case third order analysis, i.e. we obtain the required cochain at least up to and including $`\mathrm{}^2`$ terms. This will already be a substantial amount of work and is enough to expose the main phenonema. Moreover, the existence of a cochain to all orders is not really in doubt in view of the Kontsevitch universal quantisation theorem (our cochain amounts to choosing a natural ‘lifting’ of that); our results constitute a natural choice at low order and suggest that a natural choice should be possible to all further orders. We start these examples with Section 4 in which the sphere $`S^2`$ has a natural cochain $`F`$ for covariance Lie algebra $`=so(1,3)`$. The action of the Lorentz group that we use is the one on the ’sphere at infinity’ in 4-dimensional Minkowski space. We show that one obtains an associative quantisation of the sphere at least to $`O(\mathrm{}^3)`$ and that this coincides with the Fedosov quantisation to this order for the standard Levi-Civitia connection on the sphere (which is symplectic). Section 5 is the main example of interest in the paper. We show that the classical enveloping algebra $`U_{\mathrm{}}(𝔤)`$ viewed as a quantisation of $`S(𝔤)=[𝔤^{}]`$ (functions on $`𝔤^{}`$ with its Kirillov-Kostant bracket) can be viewed as a module-algebra cochain twist and that this quantizes a canonical covariant preconnection in the Poisson geometry of $`𝔤^{}`$ (we show that this is in fact the only such preconnection for all simple $`𝔤`$ other than $`sl_n`$, $`n>2`$ and even there it is the natural choice). The background covariance we use is $`=𝔤<𝔤^{}`$ and we find a suitable $`F`$ as a powerseries to $`O(\mathrm{}^3)`$ and find that it is essentially unique to this level when we demand a further condition (Section 5.4) whereby $`S(𝔤^{})=[𝔤]U()`$ twists into a local version of the group coordinate algebra $`[G]`$ (see below). In effect, we require that $`F`$ implements the Campell-Baker-Hausdorf formula by conjugation in addition to its other properties. In Section 5.5 we discuss the Duflo map in this context and argue that the reduced form of $`F`$ should be the coboundary of the Duflo operator (and hence known to all orders). Although our specific universal $`F`$ is only found to $`O(\mathrm{}^3)`$ it seems likely that these various features should extend and characterise it completely. This would be a topic for further work beyond our methods here. In Section 5.6 we demonstrate the theory on $`𝔤=>=b_+`$ the solvable Lie algebra in 2-dimensions. A version of its enveloping algebra has been proposed as ’noncommutative spacetime’\[MR\] and we exhibit a (non-unique) $`F`$ explicitly to $`O(\mathrm{}^4)`$ in this case. Section 6 completes our trio of conventional examples with the Mackey quantisation $`C^{\mathrm{}}(N)>U_{\mathrm{}}(g)`$ as a cochain quantisation of $`C^{\mathrm{}}(N)S(g)C^{\mathrm{}}(N\times g^{})`$. This extends the model in Section 5 but we need an extended cocycle and covariance Lie algebra $`=g<g^{}𝔤`$ in order to achieve this. Section 6.3 includes the case $`C^{\mathrm{}}(G)>U_{\mathrm{}}(g)`$ as a quantisation of $`T^{}G=G\times g^{}`$, where $`g`$ is the Lie algebra of a Lie group $`G`$. We follow these with the more technical example Section 7 from quantum group theory, which is simply Drinfeld’s theory for quantum groups $`_q[G]`$ reworked as a cochain twist. Here $`=gg^{\mathrm{op}}`$ acting from the left and right and $`F,\mathrm{\Phi }`$ are built from Drinfeld’s ones relating to the KZ-equations. This example is not fundamentally new but provides the role model for our view of the more conventional quantisations in the paper, so is included for completeness. Section 8 turns to the hidden nonassociativity that we have identified in the associative quantum algebras above. The most important of the many implications resulting from the cochain twist is in Section 8.2, namely the corresponding differential calculi. Because in our examples $`F`$ is not a cocycle, the exterior algebra obtained likewise by twisting is not necessarily associative, and we show that indeed it is not for our various examples. We describe the nonassociative differentials for each of our examples to order $`O(\mathrm{}^2)`$. One example is a more covariant but nonassociative differential calculus for the non-commutative spacetime in \[MR\]. Section 8.3 shows how the same philosophy can be used to construct Dirac operators in the sense of generalised ‘spectral triples’. The slight generalisation beyond the axioms in \[C\] reflects the nonassociativity. We show that such deformations are isospectral, a point of view consistent with other approaches such as \[DLSSV\]. It is also true that under the cochain quantisation scheme the original covariance becomes a (quasi)quantum group $`H_F`$ covariance, which we describe to $`\mathrm{}^2`$ in Section 8.1 for each of our examples. In the case of $`U_{\mathrm{}}(g)`$ it appears by accident to be a usual (not quasi) quantum group and to be a local version of the quantum double $`D(U(g))=U(g)<[G]`$, which is known to be a covariance quantum group of $`U(g)`$. When $`U_{\mathrm{}}(su_2)`$ is viewed as noncommutative $`^3`$ (the so-called universal fuzzy-sphere), for example, the quantum double plays the role of quantum Euclidean group\[BaMa\] motivated from 2+1 quantum gravity. In this case the curvature of the canonical preconnection or the fact that $`F`$ cannot be taken to be a cocycle represents an anomaly in this quantisation of $`^3`$. The associativity obstruction in this case can in fact be resolved by adjoining an extra ‘time’ variable and has been proposed \[M4\] as an origin of time in noncommutative differential geometry. Let us say finally that we work in a deformation-theoretic setting with all deformed expressions given by power-series in a parameter $`\mathrm{}`$ and otherwise over $``$; all constructions can be formulated more (co)algebraically over any field using comodules which would, however, be less familiar to most readers. In the main ‘examples’ sections we work only to lower degrees in $`\mathrm{}`$ for which purposes one may regard $`\mathrm{}`$ as a real parameter with the deformed product of smooth functions assumed to have these first terms in an expansion. The authors would like to thank F.W. Clarke for his assistance with some of the MATHEMATICA calculations underlying the paper and Y. Bazlov for drawing our attention to the Duflo map. ## 2. Preliminaries: module-algebra cochain twists We begin with some well known algebraic constructions, see for example the text \[M2\]. We will only need here the classical case $`H=U()`$ where $``$ is a Lie algebra and $`A=C^{\mathrm{}}(M)`$ where $`M`$ is a manifold, as the basis of the quantisation method. In that sense we use quantum group methods but the reader does not really need to know quantum group theory in any detail. This approach to quantisation as been recently used in \[AM2, M3\] for except that in the present paper the quantisation remains associative. Given a Hopf algebra $`(H,S,\mathrm{\Delta },ϵ)`$ and an invertible $`FHH`$ with $`(ϵ\mathrm{id})F=(\mathrm{id}ϵ)F=1`$, we can define a quasi-Hopf algebra $`H_F=(H,\varphi ,S_F,\mathrm{\Delta }_F,ϵ,\alpha _F,\beta _F)`$, with the same algebra and counit as $`H`$, by \[D2\] (1) $`\mathrm{\Delta }_Fh=F.\mathrm{\Delta }h.F^1,\varphi =(1F).(\mathrm{id}\mathrm{\Delta })F.(\mathrm{\Delta }\mathrm{id})F^1.(F1)^1,`$ (2) $`S_F=S,\alpha _F=(SF^{(1)}).F^{(2)},\beta _F=F^{(1)}SF^{(2)}.`$ In addition if there is a quasitriangular structure $``$ for $`H`$, then $`_F=F_{21}F^1`$ for $`H_F`$. We will call such an $`FHH`$ a 2-cochain in general, and a 2-cocycle if $`\varphi =111`$. The significance of the twisting construction is \[M1\] that it corresponds to an equivalence of categories. Thus, the category $`{}_{H_F}{}^{}`$ of left modules over $`H_F`$ is a monoidal category with tensor product operation $`^F`$ is defined using $`\mathrm{\Delta }_F`$. If $`\varphi =\varphi ^{(1)}\varphi ^{(2)}\varphi ^{(3)}`$, the associator in the category is $`\mathrm{\Phi }_{VWZ}((vw)z)=\varphi ^{(1)}v(\varphi ^{(2)}w\varphi ^{(3)}z)`$. On the other hand, this category is equivalent to the category $`{}_{H}{}^{}`$ of left modules over $`H`$ via the functor $`𝒯:{}_{H}{}^{}{}_{H}{}^{}_{F}^{}`$ which is just the identity on left $`H`$ modules and on morhphism. A monoidal functor also comes by definition with a natural transformation $`\vartheta :𝒯V^F𝒯W𝒯(VW)`$, given here by $`\vartheta (𝒯(v)^F𝒯(w))=𝒯(F^{(1)}vF^{(2)}w)`$ \[M1\]. In this way, twisting the Hopf algebra by $`F`$ deforms the entire category of modules and as such deforms any and all constructions in the category. This is the systematic ’twisting approach’ to deformation quantisation that we use. In particular, consider an algebra $`A{}_{H}{}^{}`$. This includes the requirement that multiplication $`:AAA`$ is a morphism in the category, i.e the product is $`H`$-covariant (or $`A`$ is an $`H`$-module algebra). Applying the above functor $`𝒯`$ immediately deforms the algebra to the same vector space $`A_F=A`$ and the product as a map $`𝒯(AA)𝒯A`$. Using the above natural transformation this implies a deformed product map making an algebra $`A_F{}_{H}{}^{}_{F}^{}`$ with multiplication $`ab=(F^{(1)}a)(F^{(2)}b)`$, and this is associative in the category as the image of the associativity law in the undeformed category. This module algebra cochain quantisation method was introduced in \[DGM\] and related papers at the time. Examples in the cocycle case also abound, e.g. \[MO\], but the cocycle case is not what is of interest in the present paper since in this case the associator $`\varphi `$ is trivial. Neither case of ‘module algebra twist’ should be confused with Drinfeld’s twist $`H_F`$ of the Hopf algebra $`H`$ itself. One may go further and consider also the category $`𝒜{}_{H}{}^{}`$ of $`A`$-bimodules, which also have $`H`$-actions so that the multiplications $`AVV`$ and $`VAV`$ preserve the $`H`$-action for all $`V{}_{A}{}^{}`$, and so on. Here we deform the multiplications by $`av=(F^{(1)}a).(F^{(2)}v)`$ and $`va=(F^{(1)}v).(F^{(2)}a)`$ for all $`vV`$ and $`aA`$. Similarly, if $`\mathrm{\Omega }(A)`$ is an $`H`$-covariant differential calculus in the sense of noncommutative geometry (so there is for example an exterior derivative $`\mathrm{d}:A\mathrm{\Omega }^1(A)`$ where the latter is an $`A`$-bimodule and $`\mathrm{d}`$ obeys the Leibniz rule, etc. and all maps are morphisms in $`{}_{H}{}^{}`$) then twisting any products by the action of $`F^1`$ gives a calculus $`\mathrm{\Omega }(A_F)`$ covariant under $`H_F`$. This was used for example in \[AM2\]. To this existing theory we now add some first remarks needed for the semiclassical analysis. As mentioned, the above should be understood as extended over formal power-series in a parameter $`\mathrm{}`$ or one may continue more algebraically (using a comodule twist version of the theory). Either way, we suppose that $`F^1`$ is expanded as a series (3) $`F^1`$ $`=`$ $`11+\mathrm{}G^{(1)}+\mathrm{}^2G^{(2)}+O(\mathrm{}^3).`$ This can be inverted to give (4) $`F`$ $`=`$ $`11\mathrm{}G^{(1)}+\mathrm{}^2((G^{(1)})^2G^{(2)})+O(\mathrm{}^3).`$ We can then compute $`\varphi `$ $`=`$ $`(1F).(\mathrm{id}\mathrm{\Delta })F.(\mathrm{\Delta }\mathrm{id})F^1.(F1)^1`$ $`=`$ $`111+\mathrm{}[(\mathrm{\Delta }\mathrm{id})G^{(1)}+G^{(1)}11G^{(1)}(\mathrm{id}\mathrm{\Delta })G^{(1)}]`$ $`+\mathrm{}^2[(\mathrm{\Delta }\mathrm{id})G^{(2)}+G^{(2)}1+1((G^{(1)})^2G^{(2)})+(\mathrm{id}\mathrm{\Delta })((G^{(1)})^2G^{(2)})`$ $`+(\mathrm{\Delta }\mathrm{id})G^{(1)}.(G^{(1)}1)(\mathrm{id}\mathrm{\Delta })G^{(1)}.(G^{(1)}1)(\mathrm{id}\mathrm{\Delta })G^{(1)}.(\mathrm{\Delta }\mathrm{id})G^{(1)}`$ $`(1G^{(1)}).(G^{(1)}1)(1G^{(1)}).(\mathrm{\Delta }\mathrm{id})G^{(1)}+(1G^{(1)}).(\mathrm{id}\mathrm{\Delta })G^{(1)}]+O(\mathrm{}^3).`$ Now consider a special case, where $`G^{(1)}=XY`$ (we will later suppress the summation sign), $`\mathrm{\Delta }X=1X+X1`$ and $`\mathrm{\Delta }Y=1Y+Y1`$. Then the order $`\mathrm{}`$ part of $`\varphi `$ is $`1XY+X1Y+XY11XYXY1X1Y=\mathrm{\hspace{0.17em}0}.`$ The contribution from $`G^{(1)}`$ to the order $`\mathrm{}^2`$ part of $`\varphi `$, using tildes to distinguish different copies of $`G^{(1)}`$, is $`1X\stackrel{~}{X}Y\stackrel{~}{Y}+X\stackrel{~}{X}Y\stackrel{~}{Y}1+X\stackrel{~}{X}1Y\stackrel{~}{Y}+X\stackrel{~}{X}Y\stackrel{~}{Y}+X\stackrel{~}{X}\stackrel{~}{Y}Y`$ $`+X\stackrel{~}{X}\stackrel{~}{Y}Y+\stackrel{~}{X}X\stackrel{~}{Y}YX\stackrel{~}{X}\stackrel{~}{Y}YX\stackrel{~}{X}Y\stackrel{~}{Y}1`$ $`X\stackrel{~}{X}Y\stackrel{~}{Y}X\stackrel{~}{X}1Y\stackrel{~}{Y}XY\stackrel{~}{X}\stackrel{~}{Y}X\stackrel{~}{X}Y\stackrel{~}{Y}\stackrel{~}{X}X\stackrel{~}{Y}Y`$ $`\stackrel{~}{X}XY\stackrel{~}{Y}1X\stackrel{~}{X}Y\stackrel{~}{Y}+\stackrel{~}{X}X\stackrel{~}{Y}Y+\stackrel{~}{X}XY\stackrel{~}{Y}`$ $`=`$ $`X\stackrel{~}{X}\stackrel{~}{Y}Y+\stackrel{~}{X}X\stackrel{~}{Y}YXY\stackrel{~}{X}\stackrel{~}{Y}X\stackrel{~}{X}Y\stackrel{~}{Y}`$ $`=`$ $`X\stackrel{~}{X}\stackrel{~}{Y}Y+\stackrel{~}{X}X\stackrel{~}{Y}Y\stackrel{~}{X}\stackrel{~}{Y}XYX\stackrel{~}{X}Y\stackrel{~}{Y}`$ $`=`$ $`X\stackrel{~}{X}\stackrel{~}{Y}Y+\stackrel{~}{X}[X,\stackrel{~}{Y}]YX\stackrel{~}{X}Y\stackrel{~}{Y}.`$ To summarise, if we put, for $`PHH`$, $`P`$ $`=`$ $`1P(\mathrm{\Delta }\mathrm{id})P+(\mathrm{id}\mathrm{\Delta })PP1.`$ then the expression for $`\varphi `$ becomes $`\varphi =\mathrm{\hspace{0.17em}1}11+\mathrm{}^2(X\stackrel{~}{X}\stackrel{~}{Y}Y+\stackrel{~}{X}[X,\stackrel{~}{Y}]YX\stackrel{~}{X}Y\stackrel{~}{Y}G^{(2)})+O(\mathrm{}^3).`$ Next note that, if $`A,B,C,DH`$ are also primitive (have linear coproducts like $`X,Y`$) then $`(ABCD)`$ $`=`$ $`ABCD+ABDCABCDBACD.`$ Hence if we were to put $`G^{(2)}=(G^{(1)})^2/2`$, then (5) $`\varphi =\mathrm{\hspace{0.17em}1}11+\mathrm{}^2([X,\stackrel{~}{X}]\stackrel{~}{Y}Y+2\stackrel{~}{X}[X,\stackrel{~}{Y}]YX\stackrel{~}{X}[Y,\stackrel{~}{Y}])/2+O(\mathrm{}^3).`$ The set of primitive elements in $`H`$ form a Lie algebra in $`H`$ which we can view as a Lie bialgebra with zero Lie cobracket. After twisting these elements acquire a Lie cobracket $$\delta Z=[X,Z]Y+X[Z,Y][Y,Z]XY[X,Z]$$ for all $`Z`$ in the Lie algebra and together with the $`\mathrm{}^2`$ part of $`\varphi `$ form a quasi-Lie bialgebra. This is the infinitesimal object associated to the quasi-quantum group $`H_F`$. ## 3. Poisson-compatible connections from cochain twists at the semiclassical level We start by briefly recalling the main ideas of \[BM1\]. As is well known, if one considers the quantisation of the functions $`C^{\mathrm{}}(M)`$ of a classical manifold $`M`$, the initial data usually specified is a Poisson structure defined by a bivector $`\omega `$ (in the symplectic case this is invertible with inverse (also denoted $`\omega `$) a closed 2-form). Any flat deformation-quantisation $`A_{\mathrm{}}`$ corresponds on looking at the leading part of its commutator (6) $$abba=\mathrm{}\{a,b\}+O(\mathrm{}^2)$$ to a Poisson bracket $`\{a,b\}=\omega (\mathrm{d}a,\mathrm{d}b)`$. More recently, we considered the same question for a noncommutative differential calculus $`\mathrm{\Omega }(A_{\mathrm{}})`$ quantizing the usual exterior algebra $`\mathrm{\Omega }(M)`$ but in a slightly weaker than usual setting (without assuming associativity of products involving differential forms). We found that the initial data for this at least in the symplectic case was a compatible connection $``$ defined by (7) $$a\mathrm{d}b\mathrm{d}ba=\mathrm{}_{\widehat{a}}\mathrm{d}b+O(\mathrm{}^2).$$ Here $`\widehat{a}`$ denotes the Hamiltonian vector field $`\omega (\mathrm{d}a,)=\{a,\}`$. Here $``$ is not necessarily well-defined in the Poisson case even along Hamiltonian vector fields; it could be called a partial connection where defined or we should speak more precisely of a ’preconnection’ $`\widehat{}`$ defined almost identically by (8) $$a\mathrm{d}b\mathrm{d}ba=\mathrm{}\widehat{}_a\mathrm{d}b+O(\mathrm{}^2)$$ (here $`\widehat{}`$ was called $`\gamma `$ in \[BM1\]). (In fact there is a more general notion of ‘contravariant connection’ which can also be used here, see\[H\]). The (Poisson)-compatibility condition is (9) $$_{\widehat{a}}\mathrm{d}b_{\widehat{b}}\mathrm{d}a=\mathrm{d}\{a,b\}$$ and under some mild conditions in the symplectic case becomes \[BM1\] that $``$ is a torsion free symplectic connection in the usual sense. Finally, the curvature and torsion of the connection (10) $$R_{}(\widehat{a},\widehat{b})=[_{\widehat{a}},_{\widehat{b}}]_{[\widehat{a},\widehat{b}]},T_{}(\widehat{a},\widehat{b})=_{\widehat{a}}\widehat{b}_{\widehat{b}}\widehat{a}[\widehat{a},\widehat{b}]$$ are defined in the usual way as for any conneciton. In terms of a preconnection the equations are more precisely (11) $$\widehat{}_a\mathrm{d}b\widehat{}_b\mathrm{d}a=\mathrm{d}\{a,b\},R_\widehat{}(a,b)=[\widehat{}_a,\widehat{}_b]\widehat{}_{\{a,b\}},T_\widehat{}(a,b)=\widehat{}_a\widehat{b}\widehat{}_b\widehat{a}[\widehat{a},\widehat{b}]$$ where the last term in the curvature is in view of $`[\widehat{a},\widehat{b}]=\widehat{\{a,b\}}`$. It was shown in \[BM1\] that the curvature coincides with the Jacobiator or obstruction to associativity for the differential calculus $`\mathrm{\Omega }(A_{\mathrm{}})`$ at the relevant lowest order. This was also observed recently in \[H\] where it was shown that the associative case corresponds to a zero-curvature (contravariant) connection, although we were not aware of this at the time of \[BM1\]. From a geometrical point of view, however, it would seem at the semiclassical level quite reasonable to consider symplectic or Poisson manifolds equipped with connections with curvature, which in quantisation terms would mean by \[BM1\] associative quantum algebras with nonassociative differential calculi. We have seen in Section 2 a general method to construct examples of such hybrid situations by means of cochain twists. Now we see what that amounts to at the semiclassical level. ### 3.1. Inducing connections by twisting We consider the diffeomorphism group acting on the functions on a manifold $`M`$. This action extends to the vector fields and forms on the manifold, and infinitesimally the action is called the Lie derivative. A vector field $`X`$ (i.e. an element of the Lie algebra of the diffeomorphism group) acts on forms by $`_X\xi =\varpi _X(\mathrm{d}\xi )+\mathrm{d}(\varpi _X\xi )`$, where $`\varpi _X`$ is the interior product. If $`\xi =\xi _{k_1\mathrm{}k_n}\mathrm{d}x^{k_1}\mathrm{}\mathrm{d}x^{k_n}`$, then $`_X\xi `$ $`=`$ $`X^p\xi _{k_1\mathrm{}k_n,p}\mathrm{d}x^{k_1}\mathrm{}\mathrm{d}x^{k_n}+\xi _{k_1\mathrm{}k_n}\mathrm{d}(\varpi _X(\mathrm{d}x^{k_1}\mathrm{}\mathrm{d}x^{k_n}))`$ $`=`$ $`\left(X^j\xi _{k_1\mathrm{}k_n,j}+X_{,k_i}^j\xi _{k_1\mathrm{}k_{i1},j,k_{i+1}\mathrm{}k_n}\right)\mathrm{d}x^{k_1}\mathrm{}\mathrm{d}x^{k_n}`$ We , can therefore apply the ideas of Section 2 with $`A=C^{\mathrm{}}(M)`$ and $`H=U(\mathrm{diff}(M))`$ acting via the Lie derivative on all tensorial objects. We will proceed with $`F,F^1`$ power-series having values in $`U(\mathrm{diff}(M))U(\mathrm{diff}(M))`$, however we are interested in this section only in the differential geometry resulting from the semiclassical part and not in the formal construction of these objects. We assume an expansion of the 2-cochain $`F^1=\mathrm{id}^2+\mathrm{}G^{(1)}+O(\mathrm{}^2)`$, where $`G^{(1)}=XY\mathrm{diff}(M)\mathrm{diff}(M)`$, acts on $`\mathrm{\Omega }^n(M)\mathrm{\Omega }^m(M)`$ by $`\xi \eta \xi \eta +\mathrm{}_X\xi _Y\eta +O(\mathrm{}^2)`$. Then the commutator of two functions $`a,bC^{\mathrm{}}(M)`$ is (12) $`\{a,b\}`$ $`=`$ $`\mathrm{}{\displaystyle \left((_Xa)(_Yb)(_Ya)(_Xb)\right)}+O(\mathrm{}^2)`$ (13) $`=`$ $`\mathrm{}{\displaystyle (X^iY^jX^jY^i)a_{,i}b_{,j}}+O(\mathrm{}^2).`$ Now we set $`\omega =(X_MYY_MX)\mathrm{diff}(M)_M\mathrm{diff}(M)`$ (putting summation subscripts on $`X`$ and $`Y`$ would only be confusing). This is antisymmetric, so we have $`\omega \mathrm{diff}(M)\mathrm{diff}(M)`$. We shall assume for convenience that we are in the symplectic case, where $`\omega `$ is invertible, and its inverse is a closed 2-form, otherwise more generally we assume that $`\omega `$ is a Poisson bivector, i.e. induces a Poisson bracket. A sufficient condition for the latter is that $`\varphi =111`$ when projected from $``$ to $`_M`$. Another sufficient condition is that $`G^{(1)}`$ obeys the Classical Yang-Baxter equations, but neither is necessary and we do not assume them. Now apply the same $`F^1`$ to deform the products of functions and 1-forms as explained in Section 2. This implies a connection $``$ resulting from the commutator of a function and a 1-form: $`[a,\xi _i\mathrm{d}x^i]`$ $`=`$ $`\mathrm{}{\displaystyle }X^ka_{,k}(\xi _{i,j}Y^j\mathrm{d}x^i+\xi _iY_{,j}^i\mathrm{d}x^j)(XY)+O(\mathrm{}^2)`$ $`=`$ $`\mathrm{}{\displaystyle (X^kY^jX^jY^k)a_{,k}\xi _{i,j}\mathrm{d}x^i}+{\displaystyle (X^kY_{,j}^iY^kX_{,j}^i)\xi _ia_{,k}\mathrm{d}x^j}+O(\mathrm{}^2)`$ $`=`$ $`\mathrm{}\omega ^{kj}a_{,k}\xi _{i,j}\mathrm{d}x^i+{\displaystyle (X^kY_{,j}^iY^kX_{,j}^i)\xi _ia_{,k}\mathrm{d}x^j}+O(\mathrm{}^2)`$ $`=`$ $`\mathrm{}\omega ^{kj}a_{,k}\left(\xi _{i,j}\mathrm{d}x^i+{\displaystyle \omega _{js}(X^sY_{,p}^iY^sX_{,p}^i)\xi _i\mathrm{d}x^p}\right)+O(\mathrm{}^2)`$ $`=`$ $`\mathrm{}\omega ^{kj}a_{,k}\left(\xi _{i,j}\mathrm{d}x^i\mathrm{\Gamma }_{jp}^i\xi _i\mathrm{d}x^p\right)+O(\mathrm{}^2).`$ Then the Christoffel symbols of the connection can be seen to be (14) $`\mathrm{\Gamma }_{jp}^i`$ $`=`$ $`{\displaystyle \omega _{js}(X^sY_{,p}^iY^sX_{,p}^i)}.`$ ###### Proposition 3.1.1. The connection is characterised by the equation, for $`aC^{\mathrm{}}(M)`$ and $`\xi \mathrm{\Omega }^1M`$: $$_{\widehat{a}}\xi =X(a)(\mathrm{d}Y,\xi +\varpi _Y\mathrm{d}\xi )Y(a)(\mathrm{d}X,\xi +\varpi _X\mathrm{d}\xi ).$$ Proof: Just substitute from (14) for the Christoffel symbols.$`\mathrm{}`$ The condition that $`\omega `$ is closed is $`\omega _{ir}\omega _{,k}^{rs}\omega _{sj}\mathrm{d}x^k\mathrm{d}x^i\mathrm{d}x^j=\mathrm{\hspace{0.17em}0}.`$ From \[BM1\] the connection is necessarily comapatible with the differential structure in the form (15) $`{\displaystyle \frac{\omega ^{ij}}{x^p}}`$ $`=`$ $`\omega ^{jq}\mathrm{\Gamma }_{qp}^i\omega ^{iq}\mathrm{\Gamma }_{qp}^j.`$ For the Christoffel symbols in (14) we have $`\omega ^{jq}\mathrm{\Gamma }_{qp}^i\omega ^{iq}\mathrm{\Gamma }_{qp}^j`$ $`=`$ $`{\displaystyle (X^iY_{,p}^jY^iX_{,p}^j)}{\displaystyle (X^jY_{,p}^iY^jX_{,p}^i)}`$ $`=`$ $`{\displaystyle (X^iY^jX^jY^i)_{,p}}=\omega _{,p}^{ij},`$ as expected. In summary, any manifold may potentially by quantised by choosing a cochain $`F`$ with values in $`U(\mathrm{diff}M)U(\mathrm{diff}M)[[\mathrm{}]]`$. The leading order of $`F`$ which we have denoted $`XY`$ will induce a bivector which will not necessarily be a Poisson bivector (the quantised algebra may not necessarily be associative). However, if it is, we will also have induced a Poisson-compatible connection (or more precisely a preconnection). Conversely, given, say, a symplectic manifold $`M`$ equipped with (as in Fedosov theory) a symplectic connection $``$ we can look for a suitable $`F`$ inducing these intial data to lowest order and such that $`\varphi =111`$ when tensorised over $`_M`$. This provides an alternative and more categorical approach to the quantisation problem in the spirit of Fedosov theory but now having the merit of also quantising differential calculi and all other covariant constructions, albeit with potential nonassociativity. ### 3.2. The inverse problem for $`^{2n}`$ The ‘inverse problem’ of finding $`F`$ even to lowest order (i.e. $`XY`$) such that a given Poisson bi-vector $`\omega `$ is obtained and a given symplectic or Poisson-compatible (pre)connection is obtained as above appears to be a tricky one. Here we will look at what is involved in the simplest possible case of $`^{2n}`$. We take the standard symplectic structure and note that torsion free symplectic connections for it are in 1-1 correspondence with totally symmetric Christoffel symbols $`\mathrm{\Gamma }_{abc}=\omega _{ad}\mathrm{\Gamma }_{bc}^d`$ \[GRS\]. We can easily make the canonical symplectic form for $`^{2n}`$ by adding $`\mathrm{}XY`$ terms to $`F^1`$ where $`X,Y`$ are constant vector fields. By (14) these will give zero Christoffel symbols. But then we can add further terms to $`F^1`$ of the form $`XY=f.UVUf.V`$, where $`f`$ is a real valued function and $`U,V`$ are constant vectors. Then $`\mathrm{\Gamma }_{ijp}`$ $`=`$ $`{\displaystyle \omega _{ik}\omega _{js}f_{,p}(U^sV^kV^sU^k)}.`$ If we set $`f`$ to be the linear function $`\omega _{pq}W^qx^p`$, then $`\mathrm{\Gamma }_{ijp}`$ $`=`$ $`\omega _{ik}\omega _{js}\omega _{pq}{\displaystyle W^q(U^sV^k+V^sU^k)}.`$ By adding terms of this form we can recreate any symplectic connection with constant Christoffel symbols by this form of $`F`$. Then the curvature is given by $`R_{bcd}^a`$ $`=`$ $`\omega ^{me}\omega ^{ag}(\mathrm{\Gamma }_{edb}\mathrm{\Gamma }_{gcm}\mathrm{\Gamma }_{ecb}\mathrm{\Gamma }_{gdm})`$ $`=`$ $`\omega ^{me}\omega ^{ag}\omega _{ek}\omega _{bq}\omega _{gp}\omega _{mr}(\omega _{ds}\omega _{cv}{\displaystyle }W^q(U^sV^k+V^sU^k)\stackrel{~}{W}^r(\stackrel{~}{U}^v\stackrel{~}{V}^p+\stackrel{~}{V}^v\stackrel{~}{U}^p)`$ $`\omega _{cs}\omega _{dv}{\displaystyle }W^q(U^sV^k+V^sU^k)\stackrel{~}{W}^r(\stackrel{~}{U}^v\stackrel{~}{V}^p+\stackrel{~}{V}^v\stackrel{~}{U}^p))`$ $`=`$ $`\omega _{bq}\omega _{kr}(\omega _{ds}\omega _{cv}\omega _{cs}\omega _{dv}){\displaystyle W^q(U^sV^k+V^sU^k)\stackrel{~}{W}^r(\stackrel{~}{U}^v\stackrel{~}{V}^a+\stackrel{~}{V}^v\stackrel{~}{U}^a)}`$ $`=`$ $`\omega _{bq}\omega _{kr}(\omega _{ds}\omega _{cv}\omega _{cs}\omega _{dv}){\displaystyle W^q(U^sV^k\stackrel{~}{V}^v\stackrel{~}{U}^a+V^sU^k\stackrel{~}{U}^v\stackrel{~}{V}^a)\stackrel{~}{W}^r},`$ where the tildes denote a second set of triples $`(U,V,W)`$ and the sum is over both sets. ## 4. Quantising $`S^2`$ by cochain twist Here we describe a simple example of the cochain quantisation method in Section 3. The covariance used for the twisting will be the Lorentz group and its action on $`S^2`$, a nonlinear one related to spacetime physics (the sphere at infinity in Minkowski space). This induces a quantisation not related as far as we know to the representation theoretic coadjoint orbit examples given later. ### 4.1. Some nice vector fields on the 2-sphere Our goal is to show how a natural covariance, cochain and hence quantisation arise in a nice way from the geometry of $`S^2=\{(x,y,z)^3:x^2+y^2+z^2=1\}`$. We use the standard inner product on $`^3`$. Thus, given $`\underset{¯}{v}^3`$, we have the natural vector field $`X[\underset{¯}{v}]`$ defined by $`X[\underset{¯}{v}](\underset{¯}{r})=\underset{¯}{v}\underset{¯}{r}\underset{¯}{v},\underset{¯}{r}`$ (for $`\underset{¯}{r}S^2`$) which tangent to the sphere at $`\underset{¯}{r}`$. Also at each such point we have the orbital angular moment vector field $`Y[\underset{¯}{v}](\underset{¯}{r})=\underset{¯}{v}\times \underset{¯}{r}`$, where $`\times `$ is the vector cross product. These vector fields are clearly well behaved under rotation; consider an orthogonal transformation $`TO_3()`$. Then $`T(X[\underset{¯}{v}])(\underset{¯}{r})`$ is by definition $`T(X[\underset{¯}{v}](T^1\underset{¯}{r}))=T\underset{¯}{v}T(T^1\underset{¯}{r})\underset{¯}{v},T^1\underset{¯}{r}=T\underset{¯}{v}\underset{¯}{r}T\underset{¯}{v},\underset{¯}{r}=X[T\underset{¯}{v}](\underset{¯}{r}).`$ Also we have $`T(Y[\underset{¯}{v}])(\underset{¯}{r})`$ equal to $`T(Y[\underset{¯}{v}](T^1\underset{¯}{r}))=T(\underset{¯}{v}\times T^1\underset{¯}{r})=\mathrm{det}(T).T\underset{¯}{v}\times \underset{¯}{r}=\mathrm{det}(T).Y[T\underset{¯}{v}](\underset{¯}{r}),`$ where the determinant enters by the change in sign of the vector product under a change in orientation of $`^3`$. The Lie bracket of two vector fields is defined as usual by $`[X,Y]^i=Y_{,j}^iX^jX_{,j}^iY^j`$. Thus $`[Y[\underset{¯}{v}],Y[\underset{¯}{w}]](\underset{¯}{r})`$ $`=`$ $`\underset{¯}{w}\times (\underset{¯}{v}\times \underset{¯}{r})\underset{¯}{v}\times (\underset{¯}{w}\times \underset{¯}{r})`$ $`=`$ $`\underset{¯}{w},\underset{¯}{r}\underset{¯}{v}\underset{¯}{w},\underset{¯}{v}\underset{¯}{r}\underset{¯}{v},\underset{¯}{r}\underset{¯}{w}+\underset{¯}{v},\underset{¯}{w}\underset{¯}{r}`$ $`=`$ $`\underset{¯}{w},\underset{¯}{r}\underset{¯}{v}\underset{¯}{v},\underset{¯}{r}\underset{¯}{w},`$ $`[Y[\underset{¯}{v}],X[\underset{¯}{w}]](\underset{¯}{r})`$ $`=`$ $`\underset{¯}{r}\underset{¯}{w},Y[\underset{¯}{v}](\underset{¯}{r})Y[\underset{¯}{v}](\underset{¯}{r})\underset{¯}{w},\underset{¯}{r}\underset{¯}{v}\times X[\underset{¯}{w}](\underset{¯}{r})`$ $`=`$ $`\underset{¯}{r}\underset{¯}{w},\underset{¯}{v}\times \underset{¯}{r}\underset{¯}{v}\times \underset{¯}{r}\underset{¯}{w},\underset{¯}{r}\underset{¯}{v}\times (\underset{¯}{w}\underset{¯}{r}\underset{¯}{w},\underset{¯}{r})`$ $`=`$ $`\underset{¯}{r}\underset{¯}{w},\underset{¯}{v}\times \underset{¯}{r}\underset{¯}{v}\times \underset{¯}{w}=\underset{¯}{r}\underset{¯}{v}\times \underset{¯}{w},\underset{¯}{r}\underset{¯}{v}\times \underset{¯}{w}`$ $`=`$ $`X[\underset{¯}{v}\times \underset{¯}{w}](\underset{¯}{r}),`$ $`[X[\underset{¯}{v}],X[\underset{¯}{w}]](\underset{¯}{r})`$ $`=`$ $`\underset{¯}{r}\underset{¯}{w},X[\underset{¯}{v}](\underset{¯}{r})X[\underset{¯}{v}](\underset{¯}{r})\underset{¯}{w},\underset{¯}{r}+\underset{¯}{r}\underset{¯}{v},X[\underset{¯}{w}](\underset{¯}{r})+X[\underset{¯}{w}](\underset{¯}{r})\underset{¯}{v},\underset{¯}{r}`$ $`=`$ $`\underset{¯}{r}\underset{¯}{w},\underset{¯}{v}\underset{¯}{v}\underset{¯}{w},\underset{¯}{r}+\underset{¯}{r}\underset{¯}{v},\underset{¯}{w}+\underset{¯}{w}\underset{¯}{v},\underset{¯}{r}`$ $`+\underset{¯}{r}\underset{¯}{w},\underset{¯}{r}\underset{¯}{v},\underset{¯}{r}+\underset{¯}{r}\underset{¯}{v},\underset{¯}{r}\underset{¯}{w},\underset{¯}{r}\underset{¯}{r}\underset{¯}{v},\underset{¯}{r}\underset{¯}{w},\underset{¯}{r}\underset{¯}{r}\underset{¯}{w},\underset{¯}{r}\underset{¯}{v},\underset{¯}{r}`$ $`=`$ $`\underset{¯}{v}\underset{¯}{w},\underset{¯}{r}+\underset{¯}{w}\underset{¯}{v},\underset{¯}{r}`$ whereas $`Y[\underset{¯}{v}\times \underset{¯}{w}](\underset{¯}{r})`$ $`=`$ $`(\underset{¯}{v}\times \underset{¯}{w})\times \underset{¯}{r}=\underset{¯}{r},\underset{¯}{v}\underset{¯}{w}\underset{¯}{r},\underset{¯}{w}\underset{¯}{v}.`$ Hence we have (16) $`[Y[\underset{¯}{v}],Y[\underset{¯}{w}]]=Y[\underset{¯}{v}\times \underset{¯}{w}],[Y[\underset{¯}{v}],X[\underset{¯}{w}]]=X[\underset{¯}{v}\times \underset{¯}{w}],[X[\underset{¯}{v}],X[\underset{¯}{w}]]=Y[\underset{¯}{v}\times \underset{¯}{w}].`$ In other words, the $`Y`$ fields generate rotations and the $`X`$ generare boosts of the Lie algebra $`=so(1,3)\mathrm{diff}(S^2)`$. This action has the physical interpretation mentioned above and will be used to induce the quantisaton. ### 4.2. The first order rotation invariant 2-cochain Set $`X_1=X[(1,0,0)],X_2,=X[(0,1,0)],X_3=X[(0,0,1)],`$ $`Y_1=Y[(1,0,0)],Y_2=Y[(0,1,0)],Y_3=Y[(0,0,1)].`$ Using the matrix coefficients of $`TO_3()`$ in the standard basis, $`T(X_i)={\displaystyle \underset{j}{}}T_{ji}X_j,T(Y_i)=detT.{\displaystyle \underset{j}{}}T_{ji}Y_j,`$ so under a rotation each $`X_i`$ is sent to a linear combination of the $`X_j`$ with (the important bit) constant coefficients, not general functions on $`S^2`$, and likewise with the $`Y_i`$. Using these notations, for any $`3\times 3`$ real matrix $`\kappa `$, emphasising the fact that the tensor product is over $``$, we take the lowest order part of $`F,F^1`$ in the form $`G_\kappa ^{(1)}=\kappa _{ij}X_i\underset{}{}Y_j.`$ Then, $`(TT)G_\kappa ^{(1)}=detT.\kappa _{ij}T_{ki}X_k\underset{}{}T_{sj}Y_s=detT.T_{ki}\kappa _{ij}T_{js}^TX_k\underset{}{}Y_s=detT.G_{T\kappa T^1}^{(1)}`$ Hence if we want $`\mathrm{{\rm Y}}[\kappa ]`$ to be rotation invariant, then we would like $`T\kappa T^1=\kappa `$ for all $`TSO_3()`$; we therefore take $`\kappa `$ to be (half) the identity matrix. Of course it will still have its sign changed by orientation reversing $`TO_3()`$, but this is also true of a rotation invariant symplectic form on $`S^2`$, so that this is exactly what we want. This last consideration also excludes terms of the form $`XX`$ and $`YY`$ in our ansatz for $`G^{(1)}`$. We are therefore led by rotational considerations to $`G^{(1)}=\frac{1}{2}X_iY_i`$, which we use henceforth. We use (12) to find the corresponding Poisson structure. As it is rotation invariant, we only have to evaluate $`G^{(1)}`$ at the point $`(0,0,1)`$: $`2G^{(1)}(0,0,1)`$ $`=`$ $`X_1(0,0,1)Y_1(0,0,1)+X_2(0,0,1)Y_2(0,0,1)`$ $`=`$ $`(1,0,0)(0,1,0)+(0,1,0)(1,0,0).`$ This means that, were we to reduce to $`_{C(S^2)}`$, we would get the Poisson structure corresponding to the usual symplectic form on $`S^2`$. ### 4.3. The connection It will be convenient to choose coordinates $`(x,y)^2`$ for the hemisphere $`\{(x,y,z)^3:x^2+y^2+z^2=1\mathrm{and}z>0\}`$. Then the component vector fields are given by $`X_1(x,y)=(1x^2,xy),X_2(x,y)=(xy,1y^2),X_3(x,y)=(xz,yz),`$ $`Y_1(x,y)=(0,z),Y_2(x,y)=(z,0),Y_3(x,y)=(y,x).`$ The Poisson tensor corresponding to the unique $`G^{(1)}`$ found above is $`2\omega `$ $`=`$ $`{\displaystyle \underset{i}{}}\left(X_i\underset{C(S^2)}{}Y_iY_i\underset{C(S^2)}{}X_i\right),`$ so, numbering the coordiantes $`x^1=x`$ and $`x^2=y`$, $`\omega ^{ij}`$ is antisymmetric and $`\omega ^{12}(x,y)=z`$. Taking the inverse matrix gives $`\omega _{12}=1/z`$. From (14) the Christoffel symbols are (with summation sign supressed) $`2\mathrm{\Gamma }_{1p}^i`$ $`=`$ $`\omega _{1s}(Y_k^sX_{k,p}^iX_k^sY_{k,p}^i)=(Y_k^2X_{k,p}^iX_k^2Y_{k,p}^i)/z,`$ $`2\mathrm{\Gamma }_{2p}^i`$ $`=`$ $`\omega _{2s}(Y_k^sX_{k,p}^iX_k^sY_{k,p}^i)=(Y_k^1X_{k,p}^iX_k^1Y_{k,p}^i)/z.`$ In more detail, $`2\mathrm{\Gamma }_{1p}^1`$ $`=`$ $`(Y_k^2X_{k,p}^1X_k^2Y_{k,p}^1)/z`$ $`=`$ $`(z(1x^2)_{,p}+x(xz)_{,p}(1y^2)z_{,p}yzy_{,p})/z,`$ $`2\mathrm{\Gamma }_{1p}^2`$ $`=`$ $`(Y_k^2X_{k,p}^2X_k^2Y_{k,p}^2)/z`$ $`=`$ $`(z(xy)_{,p}+x(yz)_{,p}+xy(z)_{,p}+yzx_{,p})/z,`$ $`2\mathrm{\Gamma }_{2p}^1`$ $`=`$ $`(X_k^1Y_{k,p}^1Y_k^1X_{k,p}^1)/z`$ $`=`$ $`(xyz_{,p}+xzy_{,p}z(xy)_{,p}+y(xz)_{,p})/z,`$ $`2\mathrm{\Gamma }_{2p}^2`$ $`=`$ $`(X_k^1Y_{k,p}^2Y_k^1X_{k,p}^2)/z`$ $`=`$ $`((1x^2)(z)_{,p}xzx_{,p}z(1y^2)_{,p}+y(yz)_{,p})/z.`$ Then we get, using $`z_{,1}=x/z`$ and $`z_{,2}=y/z`$, $`2\mathrm{\Gamma }_{1p}^1`$ $`=`$ $`(zxx_{,p}x^2z_{,p}(1y^2)z_{,p}yzy_{,p})/z,`$ $`2\mathrm{\Gamma }_{11}^1`$ $`=`$ $`(zx+x^3/z+(1y^2)x/z)/z=x(z^2+x^2+1y^2)/z^2,`$ $`\mathrm{\Gamma }_{11}^1`$ $`=`$ $`x(1y^2)/z^2,`$ $`2\mathrm{\Gamma }_{12}^1`$ $`=`$ $`(x^2y/z+(1y^2)y/zyz)/z=y(x^2+1y^2z^2)/z^2,`$ $`\mathrm{\Gamma }_{12}^1`$ $`=`$ $`x^2y/z^2,`$ $`2\mathrm{\Gamma }_{1p}^2`$ $`=`$ $`(z(xy)_{,p}x(yz)_{,p}xyz_{,p}+yzx_{,p})/z=\mathrm{\hspace{0.17em}2}(zx_{,p}yxyz_{,p})/z,`$ $`\mathrm{\Gamma }_{11}^2`$ $`=`$ $`y(1y^2)/z^2,`$ $`\mathrm{\Gamma }_{12}^2`$ $`=`$ $`xy^2/z^2,`$ $`2\mathrm{\Gamma }_{2p}^1`$ $`=`$ $`(xyz_{,p}+xzy_{,p}+z(xy)_{,p}y(xz)_{,p})/z=\mathrm{\hspace{0.17em}2}(xzy_{,p}yxz_{,p})/z,`$ $`\mathrm{\Gamma }_{21}^1`$ $`=`$ $`x^2y/z^2,`$ $`\mathrm{\Gamma }_{22}^1`$ $`=`$ $`x(1x^2)/z^2,`$ $`2\mathrm{\Gamma }_{2p}^2`$ $`=`$ $`((1x^2)z_{,p}xzx_{,p}+zyy_{,p}y^2z_{,p})/z,`$ $`2\mathrm{\Gamma }_{22}^2`$ $`=`$ $`((1x^2)y/z+zy+y^3/z)/z=y(1x^2+z^2+y^2)/z^2,`$ $`\mathrm{\Gamma }_{22}^2`$ $`=`$ $`y(1x^2)/z^2,`$ $`2\mathrm{\Gamma }_{21}^2`$ $`=`$ $`((1x^2)x/zxz+y^2x/z)/z=x(1x^2z^2+y^2)/z^2,`$ $`\mathrm{\Gamma }_{21}^2`$ $`=`$ $`xy^2/z^2.`$ From this we see that the connection is torsion free, and since it is also compatible with the differential structure, by \[BM1, Sec. 3\] it is also symplectic. ### 4.4. Metric compatability The metric induced from the standard embedding in $`^3`$ with the standard inner product is, in $`(x,y)`$ coordinates, $`g=\{g_{ij}\}`$ $`=`$ $`{\displaystyle \frac{1}{z^2}}\left(\begin{array}{cc}1y^2& xy\\ xy& 1x^2\end{array}\right).`$ If we organise the Christoffel symbols into the matrices $`N_k`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{\Gamma }_{k1}^1& \mathrm{\Gamma }_{k1}^2\\ \mathrm{\Gamma }_{k2}^1& \mathrm{\Gamma }_{k2}^2\end{array}\right),`$ then, using matrix multiplication, $`_kg_{ij}`$ $`=`$ $`{\displaystyle \frac{g_{ij}}{x^k}}N_k.g_{ij}g_{ij}.N_k^T.`$ In our case $`N_1={\displaystyle \frac{1}{z^2}}\left(\begin{array}{cc}x(1y^2)& y(1y^2)\\ x^2y& xy^2\end{array}\right),N_2={\displaystyle \frac{1}{z^2}}\left(\begin{array}{cc}x^2y& xy^2\\ x(1x^2)& y(1x^2)\end{array}\right),`$ and using this it can quickly be verified that the covariant derivatives of the metric vanish. Since the induced connection above was torsion free, it must be the usual Levi-Civita connection on $`S^2`$. ### 4.5. The curvature In a coordinate frame, the curvature is given by $$R_{ijk}^l=\frac{\mathrm{\Gamma }_{ki}^l}{x^j}\frac{\mathrm{\Gamma }_{ji}^l}{x^k}+\mathrm{\Gamma }_{ki}^m\mathrm{\Gamma }_{jm}^l\mathrm{\Gamma }_{ji}^m\mathrm{\Gamma }_{km}^l.$$ At the point $`x=y=0`$ we find that, to first order in $`x`$ and $`y`$, $`\mathrm{\Gamma }_{11}^1=\mathrm{\Gamma }_{22}^1=x`$, $`\mathrm{\Gamma }_{11}^2=\mathrm{\Gamma }_{22}^2=y`$ and all other Christoffel symbols vanish. Then, at that point, $$R_{ijk}^l=\frac{\mathrm{\Gamma }_{ki}^l}{x^j}\frac{\mathrm{\Gamma }_{ji}^l}{x^k}=\delta _j^l\delta _{ki}\delta _k^l\delta _{ji}.$$ Now we set $`G^{(1)}=X_iY_i/2`$ (summed over $`i`$) as above, and $`G^{(2)}=(G^{(1)})^2/2`$, then from (5) (21) $`\varphi `$ $`=`$ $`111+\mathrm{}^2([X_i,X_j]Y_jY_i+2X_j[X_i,Y_j]Y_i`$ $`X_iX_j[Y_i,Y_j])/8+O(\mathrm{}^3)`$ (23) $`=`$ $`111+\mathrm{}^2(Y[\underset{¯}{e}_i\times \underset{¯}{e}_j]Y_jY_i2X_jX[\underset{¯}{e}_i\times \underset{¯}{e}_j]Y_i`$ $`+X_iX_jY[\underset{¯}{e}_i\times \underset{¯}{e}_j])/8+O(\mathrm{}^3),`$ where the $`\underset{¯}{e}_i`$ are the usual basis vectors. Now we can expand the summations: $`X_iX_jY[\underset{¯}{e}_i\times \underset{¯}{e}_j]`$ $`=`$ $`X_{[1]}X_{[2]}Y_{[3]}X_{[1]}X_{[3]}Y_{[2]}+X_{[2]}X_{[3]}Y_{[1]}`$ $`X_{[2]}X_{[1]}Y_{[3]}+X_{[3]}X_{[1]}Y_{[2]}X_{[3]}X_{[2]}Y_{[1]},`$ $`X_jX[\underset{¯}{e}_i\times \underset{¯}{e}_j]Y_i`$ $`=`$ $`X_{[2]}X_{[3]}Y_{[1]}X_{[1]}X_{[3]}Y_{[2]}X_{[3]}X_{[2]}Y_{[1]}`$ $`+X_{[1]}X_{[2]}Y_{[3]}+X_{[3]}X_{[1]}Y_{[2]}X_{[2]}X_{[1]}Y_{[3]},`$ and using this (21) simplifies to (24) $`\varphi `$ $`=`$ $`111+\mathrm{}^2(Y[\underset{¯}{e}_i\times \underset{¯}{e}_j]Y_jX_jX[\underset{¯}{e}_i\times \underset{¯}{e}_j])Y_i/8+O(\mathrm{}^3).`$ Again expanding the summations, and assigning a name to part of (24), $`\psi `$ $`=`$ $`(Y[\underset{¯}{e}_i\times \underset{¯}{e}_j]Y_jX_jX[\underset{¯}{e}_i\times \underset{¯}{e}_j])Y_i`$ $`=`$ $`(Y_{[3]}Y_{[2]}Y_{[2]}Y_{[3]}X_{[2]}X_{[3]}+X_{[3]}X_{[2]})Y_{[1]}`$ $`+(Y_{[1]}Y_{[3]}Y_{[3]}Y_{[1]}X_{[3]}X_{[1]}+X_{[1]}X_{[3]})Y_{[2]}`$ $`+(Y_{[2]}Y_{[1]}Y_{[1]}Y_{[2]}X_{[1]}X_{[2]}+X_{[2]}X_{[1]})Y_{[3]}`$ ###### Lemma 4.5.1. If we take $`\pi `$ to be the reduction to $`_{C^{\mathrm{}}(S^2)}`$, then $`\pi \psi =0`$. ###### Proof. This is a rather long calculation done with Mathematica. Details are omitted. $`\mathrm{}`$ This corresponds to the multiplication of functions being associative to $`O(\mathrm{}^2)`$. This critical fact justifies our choice of $`G^{(2)}=(G^{(1)})^2/2`$. Note that in this case we can use (4) to calculate $`F`$ $`=`$ $`11\mathrm{}G^{(1)}+\mathrm{}^2(G^{(1)})^2/2+O(\mathrm{}^3)`$ Note that this would be consistent with $`F^1=e^{\mathrm{}G^{(1)}}=e^{\frac{\mathrm{}}{2}X_iY_i}`$ but also note that $`so(1,3)`$ is not Abelian such an exponential form will not have $`\varphi =111`$, and indeed this is not true even at order $`\mathrm{}^2`$. However, we see that when projected over $`C^{\mathrm{}}(S^2)`$ we do have that $`\varphi `$ is effectively trivial at this order on functions; in other words the example demonstrates the hybrid set up of our paper at this order. ### 4.6. The deformed algebra Following (2), if $`G^{(1)}=X_iY_i/2`$ (summation suppressed) and $`G^{(2)}=(G^{(1)})^2/2`$, then we have $`fg`$ $`=`$ $`fg+\mathrm{}(X_if)(Y_ig)/4+\mathrm{}^2(X_iX_jf)(Y_iY_jg)/8+O(\mathrm{}^3).`$ It will be convenient to continue to use the coordiantes in 4.3, in which case, using subscripts for partial differentiation, $`(X_if)(Y_ig)`$ $`=`$ $`((1x^2)f_xxyf_y)(zg_y)+(xyf_x+(1y^2)f_y)zg_x`$ $`+(xzf_xyzf_y)(yg_x+xg_y)`$ $`=`$ $`z(f_yg_xf_xg_y).`$ With rather more work, we get the following formula: $`(x^ay^b)(x^cy^d)`$ $`=`$ $`x^{a+c}y^{b+d}+\mathrm{}zx^{a+c1}y^{b+d1}(bcad)/2`$ $`+\mathrm{}^2x^{a+c2}y^{b+d2}((bc(b1)(c1)+ad(a1)(d1)2abcd)`$ $`+y^2(aca^2c+b^2cac^2b^2c^2+2abcd+ad^2a^2d^2)`$ $`+x^2(bc^2b^2c^2+a^2d+bdb^2d+2abcda^2d^2bd^2)`$ $`+x^2y^2(a+b)(c+d)(1+a+b+c+d))/8+O(\mathrm{}^3).`$ The second order part of $`fg`$, evaluated at $`x=y=0`$, is one eighth of $`{\displaystyle \frac{^2f}{x^2}}{\displaystyle \frac{^2g}{y^2}}+{\displaystyle \frac{^2f}{y^2}}{\displaystyle \frac{^2g}{x^2}}2{\displaystyle \frac{^2f}{xy}}{\displaystyle \frac{^2g}{xy}}{\displaystyle \frac{f}{x}}{\displaystyle \frac{g}{x}}{\displaystyle \frac{f}{y}}{\displaystyle \frac{g}{y}}.`$ Unsing rotation invariance, and the fact that all the Christoffel symbols vanish at $`x=y=0`$, we see that this is $`\omega ^{ij}\omega ^{kl}(_if_{,k})(_jg_{,l})g^{ij}f_{,i}g_{,j}.`$ This quantisation can be compared the Fedosov one for $`S^2`$ with the above symplectic structure and symplectic connection. We see that the second order part is not that given by the Fedosov method, as that does not have the $`g^{ij}f_{,i}g_{,j}`$ term. Note that while Fedosov gives a prescription for a quantisation that is associative on the functions to all orders from the symplectic form and connection, this is not necessarily a unique quanitsation. However we expect that our second order term $`G^{(2)}`$ may have to be modified to allow extension to all orders in $`\mathrm{}`$ as an associative algebra, and it is not obvious that this could be done within our existing 6 dimensional subalgebra of the vector fields. ## 5. Enveloping algebras $`U(g)`$ as cochain twists As an important application of the ideas in Section 3, we consider $`M=𝔤^{}`$, the dual of a Lie algebra, equipped with its standard Kirillov-Kostant Poisson structure $`\{v,w\}=[v,w]`$. Here $`S(𝔤)=[𝔤^{}]`$ i.e. we work with polynomial functions as generated by $`v𝔤`$ viewed as linear functions on $`𝔤^{}`$. ### 5.1. The cochain to lowest order. To express these canonical data as induced by a cochain twist, we seek suitable vector fields to define $`G^{(1)}=XY`$. Some natural vector fields are $`𝔤`$ itself acting by $`\mathrm{ad}`$ as mentioned above, i.e. the vector fields for the coadjoint action on $`𝔤^{}`$ from a geometrical point of view. Then there is $`𝔤^{}`$ acting by interior product on $`S(𝔤)`$, which is to say usual differentiation on $`𝔤^{}`$. These classes of vector fields generated a sub-Lie algebra $`=𝔤<𝔤^{}\mathrm{diff}(𝔤^{})`$ that turns out to be sufficient to induce the desired quantisation. Thus we take $`H=U(𝔤<𝔤^{})=U(𝔤)<S(𝔤^{})`$ acting covariantly on $`A=S(𝔤)`$. Choose a basis $`\{e_i\}`$ in $`𝔤`$, and a dual basis $`\{e^i\}`$ in $`𝔤^{}`$, and set $`F^1`$ $`=`$ $`11+\alpha \mathrm{}e_ie^i+\beta \mathrm{}e^ie_i+O(\mathrm{}^2).`$ Then $$\omega (\mathrm{d}v,\mathrm{d}w)=(\alpha \beta )e_i(v)e^i(w)e^i(v)e_i(w)=2(\alpha \beta )[v,w]$$ so we obtain the Kirillov-Kostant bracket with $`\beta \alpha =\frac{1}{2}`$. As a first consequence: ###### Proposition 5.1.1. $`𝔤^{}`$ also has on it a canonical Poisson-compatible preconnection $$\widehat{}_v\mathrm{d}w=\frac{1}{2}\mathrm{d}[v,w]$$ where $`\widehat{v}=\mathrm{ad}_v`$ is the adjoint action viewed as a vector field (in classical differential geometry, this is a derivation on $`S(𝔤)`$). The curvature and torsion are $$R(v,w)\mathrm{d}z=\frac{1}{4}[[v,w],z],T(v,w),\mathrm{d}z=\frac{1}{2}[[v,w],z]$$ ###### Proof. This comes out of the construction by using $`F`$ to deform the differential calculus, and the leading order part $`XY`$ found above. With hindsight one may check independently using the axioms in \[BM1\] that this is indeed a canonical Poisson-compatible preconnection for the Kirillov-Kostant bracket. We compute its curvature as $$R(v,w)(\mathrm{d}z)=(\widehat{}_v\widehat{}_w\widehat{}_w\widehat{}_v_{\widehat{\{v,w\}}})(\mathrm{d}z)=\mathrm{d}([v,[w,z]][w,[v,z]]2[[v,w],z])/4$$ with the result stated in view of the Jacobi identity in the Lie algebra. Similarly $$T(v,w),\mathrm{d}z=\widehat{}_v\widehat{w}\widehat{}_w\widehat{v}[\widehat{v},\widehat{w}],\mathrm{d}z=\widehat{v}(\widehat{w}(z))\widehat{w},\widehat{}_v\mathrm{d}z(vw)[[v,w],z]$$ with the result stated again on using the Jacobi identity. Incidentally, leaving out the $`1/2`$ gives a preconnection with zero curvature, but it is not Poisson-compatible. $`\mathrm{}`$ Note that the most general translation-invariant $`\widehat{}`$ in this quantisation is given by the analysis of \[BM1\] as of the form $$\widehat{}_v\mathrm{d}w=\frac{1}{2}\mathrm{d}[v,w]+\mathrm{d}\widehat{\mathrm{\Xi }}(v,w)$$ where $`\widehat{\mathrm{\Xi }}:𝔤𝔤𝔤`$ is some symmetric linear map. This follows from regarding $`𝔤^{}`$ as an Abelian Lie group and applying the theory in \[BM1, Sec. 4.1\]. The operations $`L^{}`$ and $`R^{}`$ translating differentials back to the origin are trivial so that $`\widehat{}_v\mathrm{d}w=\mathrm{d}\mathrm{\Xi }(v,w)`$ is defined by a map $`\mathrm{\Xi }`$ with arbitrary symmetric part, which we denote $`\widehat{\mathrm{\Xi }}`$, and antisymmetric part the same as above. On the other hand, if we further demand background ‘rotational’ invariance in the sense of $`\mathrm{ad}`$-invariance under $`𝔤`$ (which becomes covariance of the calculus under the quantum double $`D(U(𝔤))`$ after quantisation) this corresponds to $`\widehat{\mathrm{\Xi }}`$ symmetric and ad-invariant. ###### Theorem 5.1.2. For all simple $`𝔤`$ other than $`sl_n`$, $`n>2`$ the canonical preconnection in Proposition 5.1.1 is the only translation and $`𝔤`$-invariant one on $`S(𝔤)=[𝔤^{}]`$. For $`𝔤=sl_n`$, $`n>2`$ there is a 1-parameter moduli space of such covariant $`\widehat{}`$ but they all have curvature. Hence for all simple $`𝔤`$ any covariant differential calculus on $`U_{\mathrm{}}(𝔤)`$ with classical dimensions is necessarily nonassociative. ###### Proof. By the same arguments from invariant theory as in the proof of \[BM1, Theorem 4.20\] (but now in a different context), basically from Kostant’s work, there is no nonzero symmetric ad-invariant map $`\widehat{\mathrm{\Xi }}:𝔤𝔤𝔤`$ for $`𝔤`$ simple other than for $`sl_n`$, $`n>2`$. Hence $`\widehat{\mathrm{\Xi }}=0`$ and $`\widehat{}`$ has to be the one in Proposition 5.1.1. For $`sl_n`$, $`n>2`$ one has the possibility of a 1-parameter family via the invariant totally symmetric trilinear form viewed as the map $`\widehat{\mathrm{\Xi }}`$. In this case $$R(v,w)\mathrm{d}z=\frac{1}{4}\mathrm{d}[[v,w],z]+\mathrm{d}\left(\widehat{\mathrm{\Xi }}(v,\widehat{\mathrm{\Xi }}(w,z))\widehat{\mathrm{\Xi }}(w,\widehat{\mathrm{\Xi }}(v,z))\right)$$ since the terms linear in $`\widehat{\mathrm{\Xi }}`$ cancel using its $`\mathrm{ad}`$-invariance. We have to show that there are always $`v,w,z`$ with the curvature expression nonzero. To do this, note that the symmetric trilinear is a cubic polynomial on $`sl_n`$ which on $`vsl_n`$ has value $`I(v,v,v)=I(v)`$, say (e.g. for $`sl_3`$ we have $`I(v)=det(v)`$). We can reconstruct the full trilinear from this by polarisation, e.g. $$I(v,w,w)=\frac{1}{6}(I(v+2w)2I(v+w)+I(v))I(w)$$ and we define $`\widehat{\mathrm{\Xi }}(v,w)=I(v,w,e_i)e_j\kappa ^{ij}`$ where $`\kappa ^{ij}`$ is the inverse matrix of the Killing form (not necessarily normalised). We fix $`v,w`$ diagonal (i.e. in the standard Cartan subalgebra of $`sl_n`$) and focus on $$R(v,w)\mathrm{d}w=I(v,e_i,w)I(w,v,e_j)\kappa ^{ij}I(w,e_i,w)I(v,v,e_j)\kappa ^{ij}.$$ We will show that this can be arranged to be non-zero. Note that if $`t`$ lies in the Cartan subalgebra and $`zsl_n`$, then $`\mathrm{ad}`$-invariance $`I([t,v],w,z)+I(v,[t,w],z)+I(v,w,[t,z])=0`$ means $`I(v,w,[t,z])=0`$. We conclude that if $`z`$ is a root vector then, $`I(v,w,z)=0`$ (since $`[t,z]`$ is a nonzero multiple of $`z`$). Hence $`I(v,w,z)`$ vanishes for all $`z`$ in the space spanned by the nonzero root vectors, which is to say the orthogonal complement of the Cartan with respect to the Killing form (it is the space of matrices in $`sl_n`$ with zero diagonal). Hence we let $`\{e_a\}`$ be a basis of the Cartan subalgebra completed to a basis of $`sl_n`$ taken from this complement. It means that we can compute $`R(v,w)\mathrm{d}w`$ using only a sum over the $`e_a,e_b`$ in place of $`e_i,e_j`$ in the expression above. For $`sl_3`$ we take $`t_1=e_{11}e_{22}`$, $`t_2=e_{22}e_{33}`$ in the Cartan. Then $`e_1=t_1`$ and $`e_2=t_1+\frac{1}{2}t_2=\frac{1}{2}(e_{11}+e_{22})e_{33}`$ are a basis with $`\kappa ^{ab}=\mathrm{diag}(1/2,2/3)`$. We also compute $$I_{111}I(t_1,t_1,t_1)=I_{222}I(t_2,t_2,t_2)=1,I_{112}I(t_1,t_1,t_2)=\frac{3}{2},I_{122}I(t_1,t_2,t_2)=\frac{5}{6}$$ using the polarisation formula above. Hence setting $`v=t_1`$, $`w=t_2`$ we compute $`R(v,w)\mathrm{d}w`$ $`=`$ $`{\displaystyle \frac{1}{2}}I(t_1,t_2,t_1)^2+{\displaystyle \frac{2}{3}}I(t_1,t_2,t_1+{\displaystyle \frac{1}{2}}t_2)^2`$ $`{\displaystyle \frac{1}{2}}I(t_1,t_1,t_1)I(t_2,t_2,t_1){\displaystyle \frac{2}{3}}I(t_1,t_1,t_1+{\displaystyle \frac{1}{2}}t_2)I(t_2,t_2,t_1+{\displaystyle \frac{1}{2}}t_2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}I_{111}^2+{\displaystyle \frac{2}{3}}(I_{112}+{\displaystyle \frac{1}{2}}I_{122})^2{\displaystyle \frac{1}{2}}I_{111}I_{122}{\displaystyle \frac{2}{3}}(I_{111}+{\displaystyle \frac{1}{2}}I_{112})(I_{122}+{\displaystyle \frac{1}{2}}I_{222})>0`$ for the values stated. This proves the result for $`sl_3`$. For $`sl_n`$ the trilinear is given by $`I(v)=_{i<j<k}v^iv^jv^k`$ in terms of the diagonal entries of $`v`$ in the Cartan. We take the same $`v=t_1,w=t_2`$ as above but viewed in the standard way inside $`sl_n`$ rather than $`sl_3`$ and the $`e_1,e_2`$ completed to a diagonal basis for $`\kappa `$. Then the computation reduces to the same one as above for $`sl_3`$. $`\mathrm{}`$ In summary, these results tells us that for simple $`𝔤`$ we are going to necessarily have to work with nonassociative differentials, and for all Lie algebras $`𝔤`$ the canonical ‘universal’ choice (which is often the only choice) at the lowest order level is the one in Proposition 5.1.1. We therefore focus on this and have seen that it is indeed given by a cochain twist at lowest order. We next want to extend Proposition 5.1.1 to find $`F,F^1`$ at least to order $`O(\mathrm{}^3)`$. To do this we first look at the product of $`U_{\mathrm{}}(𝔤)`$ on monomials. Computations have been done with MATHEMATICA. ### 5.2. The Campbell-Baker-Hausdorff product We consider which $`F`$ induce not only the above semiclassical data but the actual product of $`U_{\mathrm{}}(𝔤)`$ as a star-product quantisation of $`S(g)`$. Here $`U_{\mathrm{}}(𝔤)`$ denotes the tensor algebra on $`𝔤`$ with relations $`vwwv=\mathrm{}[v,w]`$ in terms of the Lie bracket of $`𝔤`$. This is a deformation of $`S(𝔤)`$ by the linear map linear map $`\phi :S(𝔤)U_{\mathrm{}}(𝔤)`$ given by a sum over permutations $`\phi (v_1\mathrm{}v_n)`$ $`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle \underset{\kappa S_n}{}}v_{\kappa (1)}\mathrm{}v_{\kappa (n)}.`$ As explained in \[G\], $`\phi `$ is a 1-1 correspondence, and we define a deformed multiplication $``$ on $`S(𝔤)`$ by $`\underset{¯}{v}\underset{¯}{w}=\phi ^1(\phi (\underset{¯}{v}).\phi (\underset{¯}{w}))`$. As examples, $`\phi (v).\phi (w)`$ $`=`$ $`(vw+wv)/2+\mathrm{}[v,w]/2`$ $`=`$ $`\phi (vw)+\mathrm{}\phi ([v,w])/2,`$ $`\phi (v^2).\phi (w)`$ $`=`$ $`(v^2w+vwv+wv^2)/3+\mathrm{}(2v[v,w]+[v,w]v)/3`$ $`=`$ $`\phi (v^2w)+\mathrm{}\phi (v[v,w])/2+\mathrm{}^2\phi ([v[v,w]])/6.`$ This is related to the CBH formula as follows: since $`e^v`$ in $`S(𝔤)`$ maps under $`\varphi `$ to $`e^v`$ in $`U_{\mathrm{}}(𝔤)`$ (similarly for any power series in $`v`$) we have $$e^ve^w=\varphi ^1(\varphi (e^v)\varphi (e^w))=\varphi ^1(e^ve^w)=\varphi ^1(e^{C_{\mathrm{}}(v,w)})=e^{C_{\mathrm{}}(v,w)}$$ where $`C_{\mathrm{}}(v,w)`$ is the CBH power series for the product of two exponentials in $`U(𝔤)`$ with the insertion of powers of $`\mathrm{}`$ for each commutator in the Lie algebra. ###### Lemma 5.2.1. For $`v_1\mathrm{}v_n`$ a symmetric product of elements of $`𝔤`$, and $`w𝔤`$, $`wv_1\mathrm{}v_n`$ $`=`$ $`\phi (wv_1\mathrm{}v_n)\mathrm{}n\phi ([v_1,w]v_2\mathrm{}v_n)/2`$ $`+\mathrm{}^2n(n1)\phi ([v_1,[v_2,w]]v_3\mathrm{}v_n)/12+O(\mathrm{}^3).`$ Proof: Set $`0ina_i^0`$ $`=`$ $`v_1\mathrm{}v_iw\mathrm{}v_n,`$ $`0in1a_i^1`$ $`=`$ $`v_1\mathrm{}v_i[v_{i+1},w]\mathrm{}v_n,`$ $`0in2a_i^2`$ $`=`$ $`v_1\mathrm{}v_i[v_{i+1},[v_{i+2},w]]\mathrm{}v_n.`$ Then in $`U_{\mathrm{}}(𝔤)`$ we have $`a_i^ma_{i+1}^m`$ $`=`$ $`\mathrm{}a_i^{m+1},`$ and this gives, for $`i>0`$, (25) $`a_i^m`$ $`=`$ $`a_0^m+\mathrm{}(a_0^{m+1}+\mathrm{}+a_{i1}^{m+1}).`$ From this we get (26) $`a_0^0+\mathrm{}+a_n^0`$ $`=`$ $`(n+1)a_0^0+\mathrm{}(na_0^1+(n1)a_1^1+\mathrm{}+a_{n1}^1)`$ (28) $`=`$ $`(n+1)a_0^0+\mathrm{}(n+1)(a_0^1+\mathrm{}+a_{n1}^1)/2`$ $`+\mathrm{}((n1)a_0^1+(n3)a_1^1+\mathrm{}+(1n)a_{n1}^1)/2.`$ Now we use (25) again to get $`{\displaystyle \underset{j:0jn1}{}}(n12j)a_j^1`$ $`=`$ $`\mathrm{}{\displaystyle \underset{j:0jn1}{}}(n12j){\displaystyle \underset{i:0ij1}{}}a_i^2`$ $`=`$ $`\mathrm{}{\displaystyle \underset{i:0in2}{}}a_i^2{\displaystyle \underset{j:i+1jn1}{}}(n12j)`$ $`=`$ $`\mathrm{}{\displaystyle \underset{i:0in2}{}}a_i^2(i+1)(n1i),`$ and from this (26) becomes the following, where we use (25) again to get the last equality: (30) $`a_0^0+\mathrm{}+a_n^0`$ $`=`$ $`(n+1)a_0^0+\mathrm{}(n+1)(a_0^1+\mathrm{}+a_{n1}^1)/2`$ $`{\displaystyle \frac{\mathrm{}^2}{2}}{\displaystyle \underset{i:0in2}{}}a_i^2(i+1)(n1i)`$ (32) $`=`$ $`(n+1)a_0^0+\mathrm{}(n+1)(a_0^1+\mathrm{}+a_{n1}^1)/2`$ $`{\displaystyle \frac{n(n+1)\mathrm{}^2}{12}}(a_0^2+\mathrm{}+a_{n2}^2)+O(\mathrm{}^3)`$ as required. $`\mathrm{}`$ ###### Lemma 5.2.2. For $`v_1\mathrm{}v_n`$ and $`w_0\mathrm{}w_m`$ symmetric products of elements of $`𝔤`$, $`w_0\phi (w_1\mathrm{}w_mv_1\mathrm{}v_n)`$ $`=`$ $`\phi (w_0w_1\mathrm{}w_mv_1\mathrm{}v_n)`$ $`\mathrm{}n\phi ([v_1,w_0]w_1\mathrm{}w_mv_2\mathrm{}v_n)/2`$ $`+\mathrm{}^2\phi (n(n1)[v_1,[v_2,w_0]]w_1\mathrm{}w_mv_3\mathrm{}v_n`$ $`+nm[w_1,[v_1,w_0]]w_2\mathrm{}w_mv_2\mathrm{}v_n)/12+O(\mathrm{}^3)`$ $`w_0\phi ([v_1,w_1]w_2\mathrm{}w_mv_2\mathrm{}v_n)`$ $`=`$ $`\phi ([v_1,w_0]w_1\mathrm{}w_mv_2\mathrm{}v_n)`$ $`\mathrm{}\phi ([w_0,[w_1,v_1]]w_2\mathrm{}w_mv_2\mathrm{}v_n`$ $`+(n1)[v_1,w_0][v_2,w_1]w_2\mathrm{}w_mv_3\mathrm{}v_n)/2+O(\mathrm{}^2).`$ Proof: Using 5.2.1 and being careful about counting permutations, we get, for $`u𝔤`$, to $`O(\mathrm{}^3)`$, $`u\phi (w_1\mathrm{}w_mv_1\mathrm{}v_n)`$ $`=`$ $`\phi (uw_1\mathrm{}w_mv_1\mathrm{}v_n)`$ $`\mathrm{}(n+m)n\phi ([v_1,u]w_1\mathrm{}w_mv_2\mathrm{}v_n)/(2(n+m))`$ $`\mathrm{}(n+m)m\phi ([w_1,u]w_2\mathrm{}w_mv_1\mathrm{}v_n)/(2(n+m))`$ $`+\mathrm{}^2(n+m)(n+m1)(1/12)\phi (`$ $`n(n1)[v_1,[v_2,u]]w_1\mathrm{}w_mv_3\mathrm{}v_n/((n+m)(n+m1))`$ $`+nm[v_1,[w_1,u]]w_2\mathrm{}w_mv_2\mathrm{}v_n/((n+m)(n+m1))`$ $`+nm[w_1,[v_1,u]]w_2\mathrm{}w_mv_2\mathrm{}v_n/((n+m)(n+m1))`$ $`+m(m1)[w_1,[w_2,u]]w_3\mathrm{}w_mv_1\mathrm{}v_n/((n+m)(n+m1))).`$ Putting $`u=w_0`$ and supposing that $`w_0\mathrm{}w_m`$ is symmetrised, this reduces to the first part of the statement. Next, to $`O(\mathrm{}^2)`$, $`u\phi ([v_1,w_1]w_2\mathrm{}w_mv_2\mathrm{}v_n)`$ $`=`$ $`\phi (u[v_1,w_1]w_2\mathrm{}w_mv_2\mathrm{}v_n)`$ $`\mathrm{}(n+m1)(1/2)\phi (`$ $`[[v_1,w_1],u]w_2\mathrm{}w_mv_2\mathrm{}v_n/(n+m1)`$ $`+(m1)[w_2,u][v_1,w_1]w_3\mathrm{}w_mv_2\mathrm{}v_n/(n+m1)`$ $`+(n1)[v_2,u][v_1,w_1]w_2\mathrm{}w_mv_3\mathrm{}v_n/(n+m1)).`$ If $`w_0\mathrm{}w_m`$ is symmetrised, this reduces to the second part of the statement. $`\mathrm{}`$ ###### Proposition 5.2.3. For symmetric $`w_1\mathrm{}w_m`$ and $`v_1\mathrm{}v_n`$: $`w_1\mathrm{}w_mv_1\mathrm{}v_n`$ $`=`$ $`\phi (w_1\mathrm{}w_mv_1\mathrm{}v_n)`$ $`\mathrm{}mn\phi ([v_1,w_1]w_2\mathrm{}w_mv_2\mathrm{}v_n)/2`$ $`+\mathrm{}^2n(n1)m(m1)\phi ([v_1,w_1][v_2,w_2]w_3\mathrm{}w_mv_3\mathrm{}v_n)/8`$ $`+\mathrm{}^2n(n1)m\phi ([v_1,[v_2,w_1]]w_2\mathrm{}w_mv_3\mathrm{}v_n)/12`$ $`+\mathrm{}^2m(m1)n\phi ([w_1,[w_2,v_1]]w_3\mathrm{}w_mv_2\mathrm{}v_n)/12+O(\mathrm{}^3).`$ Proof: By induction on $`m`$. We will suppose that, for fixed $`n`$ and symmetric $`w_1\mathrm{}w_m`$, $`w_1\mathrm{}w_mv_1\mathrm{}v_n`$ $`=`$ $`\phi (w_1\mathrm{}w_mv_1\mathrm{}v_n)\mathrm{}\alpha _m\phi ([v_1,w_1]w_2\mathrm{}w_mv_2\mathrm{}v_n)`$ $`+\mathrm{}^2(\beta _m\phi ([v_1,w_1][v_2,w_2]w_3\mathrm{}w_mv_3\mathrm{}v_n)`$ $`+\gamma _m\phi ([v_1,[v_2,w_1]]w_2\mathrm{}w_mv_3\mathrm{}v_n)`$ $`+\delta _m\phi ([w_1,[w_2,v_1]]w_3\mathrm{}w_mv_2\mathrm{}v_n))+O(\mathrm{}^3).`$ Using associativity of the $``$ product, $`w_0\mathrm{}w_mv_1\mathrm{}v_n`$ $`=`$ $`\phi (w_0w_1\mathrm{}w_mv_1\mathrm{}v_n)`$ $`\mathrm{}n\phi ([v_1,w_0]w_1\mathrm{}w_mv_2\mathrm{}v_n)/2`$ $`+\mathrm{}^2\phi (n(n1)[v_1,[v_2,w_0]]w_1\mathrm{}w_mv_3\mathrm{}v_n`$ $`nm[w_0,[w_1,v_1]]w_2\mathrm{}w_mv_2\mathrm{}v_n)/12`$ $`\mathrm{}\alpha _m\phi ([v_1,w_0]w_1\mathrm{}w_mv_2\mathrm{}v_n)`$ $`+\mathrm{}^2\alpha _m\phi ([w_0,[w_1,v_1]]w_2\mathrm{}w_mv_2\mathrm{}v_n`$ $`+(n1)[v_1,w_0][v_2,w_1]w_2\mathrm{}w_mv_3\mathrm{}v_n)/2`$ $`+\mathrm{}^2(\beta _m\phi ([v_1,w_0][v_2,w_1]w_2\mathrm{}w_mv_3\mathrm{}v_n)`$ $`+\gamma _m\phi ([v_1,[v_2,w_0]]w_1\mathrm{}w_mv_3\mathrm{}v_n)`$ $`+\delta _m\phi ([w_0,[w_1,v_1]]w_2\mathrm{}w_mv_2\mathrm{}v_n))+O(\mathrm{}^3).`$ This gives the recursive equations and initial conditions $`\alpha _{m+1}`$ $`=`$ $`\alpha _m+n/2,\alpha _1=n/2,`$ $`\beta _{m+1}`$ $`=`$ $`\beta _m+\alpha _m(n1)/2,\beta _1=\mathrm{\hspace{0.17em}0},`$ $`\gamma _{m+1}`$ $`=`$ $`\gamma _m+n(n1)/12,\gamma _1=n(n1)/12,`$ $`\delta _{m+1}`$ $`=`$ $`\delta _mnm/12+\alpha _m/2,\delta _1=\mathrm{\hspace{0.17em}0}.`$ From this we get $`\alpha _m=nm/2`$, $`\gamma _m=n(n1)m/12`$, $`\delta _m=m(m1)n/12`$ and $`\beta _m=n(n1)m(m1)/8`$. $`\mathrm{}`$ ### 5.3. Cochain for the deformed product of $`S(𝔤)`$ Choose a dual basis $`e_i𝔤`$ and $`e^i𝔤^{}`$. Let $`𝔤`$ act on $`S(𝔤)`$ by the adjoint, and $`𝔤^{}`$ act by evaluation. Set $`Q_1=e_ie^i`$ and $`Q_2=e^ie_i`$, and let $`\mu `$ stand for multiplication. Then for symmetric $`\underset{¯}{w}=w_1\mathrm{}w_m`$ and $`\underset{¯}{v}=v_1\mathrm{}v_n`$: $`\mu (e_ie^i)(\underset{¯}{w}\underset{¯}{v})`$ $`=`$ $`n[v_1,w_1\mathrm{}w_m]v_2\mathrm{}v_n`$ $`=`$ $`nm[v_1,w_1]w_2\mathrm{}w_mv_2\mathrm{}v_n`$ $`\mu (e^ie_i)(\underset{¯}{w}\underset{¯}{v})`$ $`=`$ $`mw_2\mathrm{}w_m[w_1,v_1\mathrm{}v_n]`$ $`=`$ $`nm[v_1,w_1]w_2\mathrm{}w_mv_2\mathrm{}v_n.`$ We have quadratic terms of the form $`Q_1^2`$, $`Q_2^2`$ and $`:Q_1Q_2:`$ (with $`::`$ denoting a normal ordering with elements of $`𝔤^{}`$ being put on the right), which are respectively $`\mu (e_ie_je^ie^j)(\underset{¯}{w}\underset{¯}{v})`$ $`=`$ $`n(n1)[v_1,[v_2,w_1\mathrm{}w_m]]v_3\mathrm{}v_n`$ $`=`$ $`n(n1)m[v_2,[v_1,w_1]w_2\mathrm{}w_m]v_3\mathrm{}v_n`$ $`=`$ $`n(n1)m[v_2,[v_1,w_1]]w_2\mathrm{}w_mv_3\mathrm{}v_n`$ $`+n(n1)m(m1)[v_1,w_1][v_2,w_2]w_3\mathrm{}w_mv_3\mathrm{}v_n,`$ $`\mu (e^ie^je_ie_j)(\underset{¯}{w}\underset{¯}{v})`$ $`=`$ $`m(m1)w_3\mathrm{}w_m[w_1,[w_2,v_1\mathrm{}v_n]]`$ $`=`$ $`nm(m1)w_3\mathrm{}w_m[w_1,[w_2,v_1]v_2\mathrm{}v_n]`$ $`=`$ $`nm(m1)w_3\mathrm{}w_m[w_1,[w_2,v_1]]v_2\mathrm{}v_n`$ $`+n(n1)m(m1)w_3\mathrm{}w_m[w_2,v_2][w_1,v_1]v_3\mathrm{}v_n,`$ $`\mu (e_ie^je_je^i)(\underset{¯}{w}\underset{¯}{v})`$ $`=`$ $`n[v_1,e^j(\underset{¯}{w})].e_j(v_2\mathrm{}v_n)`$ $`=`$ $`nm[v_1,w_2\mathrm{}w_m].[w_1,v_2\mathrm{}v_n]`$ $`=`$ $`n(n1)m(m1)[v_1,w_1][v_2,w_2]w_3\mathrm{}w_mv_3\mathrm{}v_n.`$ If we set (33) $`F^1`$ $`=`$ $`11+\mathrm{}(\alpha Q_1+\beta Q_2)+\mathrm{}^2(2Q_1^2+2Q_2^2+:Q_1Q_2:)/24+O(\mathrm{}^3),`$ where $`\alpha \beta =1/2`$, then we recover the CBH product to $`O(\mathrm{}^3)`$. Note that we could add any multiple of $`Q_1Q_2+Q_2^2`$ or $`Q_2Q_1+Q_1^2`$ to $`G^{(2)}`$ and still get the same product to $`O(\mathrm{}^3)`$. Also we have the equation $`\mu (:Q_1Q_2:^R+:Q_1Q_2:+Q_1^2+Q_2^2)(\underset{¯}{w}\underset{¯}{v})`$ $`=`$ $`nmw_2\mathrm{}w_mv_2\mathrm{}v_n.e^i([v_1,[w_1,e_i]]),`$ where $`:Q_1Q_2:^R=e^je_ie^ie_j`$ is the reversed normal order. Now $`e^i([v,[w,e_i]])`$ (summed over $`i`$) is the trace of $`\mathrm{ad}_v\mathrm{ad}_w`$, that is $`v,w`$, where $`,`$ is the killing form. If we set $`e_i,e_j=\kappa _{ij}`$, then $`\mu (:Q_1Q_2:^R+:Q_1Q_2:+Q_1^2+Q_2^2+\kappa _{ij}e^ie^j)(\underset{¯}{w}\underset{¯}{v})`$ $`=`$ $`0.`$ ### 5.4. Improved cochain for the deformed coproduct on $`S(𝔤^{})`$ Here we consider a more general covariant ansatz but show that a further requirement relating to the coproduct of $`S(𝔤)`$ again leads to a unique answer. Thus, we can write a more general expression for $`G^{(2)}`$ as (35) $`G^{(2)}`$ $`=`$ $`(2Q_1^2+2Q_2^2+:Q_1Q_2:)/24+\gamma (Q_1+Q_2)Q_1+\delta (Q_1+Q_2)Q_2`$ $`+\zeta (:Q_1Q_2:^R+:Q_1Q_2:+Q_1^2+Q_2^2+\kappa _{ij}e^ie^j),`$ where $`\gamma ,\delta `$ and $`\zeta `$ are constants, and we shall use this instead of the $`\mathrm{}^2`$ term in (33). By the above discussion, this still gives the correct deformed product on $`S(𝔤)`$. We use $`G^{(1)}=\alpha Q_1+\beta Q_2`$. Then from (4) the deformed coproduct is given by (36) $`\mathrm{\Delta }_F(x)=\mathrm{\Delta }(x)+\mathrm{}[\mathrm{\Delta }(x),G^{(1)}]+\mathrm{}^2([\mathrm{\Delta }(x),G^{(2)}]G^{(1)}[\mathrm{\Delta }(x),G^{(1)}])+O(\mathrm{}^3).`$ For $`x𝔤^{}`$ we have $`\mathrm{\Delta }(x)=x1+1x`$, and using the fact that elements of $`𝔤^{}`$ commute, we find $`[\mathrm{\Delta }(x),G^{(1)}]=\alpha [x,e_i]e^i+\beta e^i[x,e_i].`$ The coefficient of $`\mathrm{}^2`$ in (36) is $`\left(\gamma [x,e_i]e_j+(\gamma \alpha ^2)e_j[x,e_i]+({\displaystyle \frac{1}{12}}+\zeta )[[x,e_i],e_j]\right)e^ie^j`$ $`+e^ie^j\left(\delta [x,e_i]e_j+(\delta \beta ^2)e_j[x,e_i]+({\displaystyle \frac{1}{12}}+\zeta )[[x,e_i],e_j]\right)`$ $`+e^i[x,e_j]\left((\gamma +{\displaystyle \frac{1}{24}}\alpha \beta +\zeta )e_ie^j+(\delta +\zeta )e^je_i\right)`$ $`+\left((\gamma +\zeta )e^ie_j+(\delta +{\displaystyle \frac{1}{24}}\alpha \beta +\zeta )e_je^i\right)[x,e_i]e^j.`$ To ensure that this is in $`S(𝔤^{})S(𝔤^{})`$ we require that $`\gamma =\alpha ^2/2`$, $`\delta =\beta ^2/2`$ and $`(\alpha \beta )^2=4\zeta 1/12`$. We already have $`\alpha \beta =1/2`$, so we get $`\zeta =1/12`$. Now the coefficient of $`\mathrm{}^2`$ in (36) is: $`\mathrm{\Delta }_F(x)`$ $`=`$ $`x1+1x+\mathrm{}(\alpha [x,e_i]e^i+\beta e^i[x,e_i]`$ $`+\mathrm{}^2\alpha ^2[[x,e_i],e_j]e^ie^j/2+e^ie^j\beta ^2[[x,e_i],e_j]/2`$ $`+(\beta ^2/21/12)e^i[x,e_j][e^j,e_i]+(\alpha ^2/21/12)[e^i,e_j][x,e_i]e^j.`$ Also putting these values into the general form of $`F^1`$ we obtain: ###### Theorem 5.4.1. The cochain $`F^1`$ $`=`$ $`11+\mathrm{}(\alpha Q_1+\beta Q_2)`$ $`+\mathrm{}^2({\displaystyle \frac{\alpha ^2}{2}}(Q_1+Q_2)Q_1+{\displaystyle \frac{\beta ^2}{2}}(Q_1+Q_2)Q_2{\displaystyle \frac{1}{24}}(:Q_1Q_2:+2:Q_1Q_2:^R+2\kappa _{ij}e^ie^j))+O(\mathrm{}^3)`$ for $`\alpha \beta =\frac{1}{2}`$ reproduces the product of $`U_{\mathrm{}}(𝔤)`$ to the relevant order and has the property that $`\mathrm{\Delta }_F(S(𝔤^{}))S(𝔤^{})S(𝔤^{})`$ to the relevant order. It appears that the up to the choice of how $`\frac{1}{2}`$ is split between $`\alpha `$ and $`\beta `$, the cochain $`F`$ is determined at higher orders by the properties in the theorem to hold for any Lie algebra and the requirement of having a $`𝔤`$-invariant form. Although we will not give a formal proof, let us explain the underlying reason here. First we note that locally near the identity we may identify the Lie algebra with the connected and simply connected Lie group $`G`$ associated to it, i.e. $`S(𝔤^{})_F_{loc}[G]`$, where the latter denotes functions defined near the identity. This is by $$\mathrm{\Theta }:S(𝔤^{})𝔤^{}xf_x(e^\mathrm{}v)=\mathrm{}^1x,v$$ (i.e. the generators appear as logarithmic coordinates on the Lie group). Let us now show that $$((\mathrm{\Theta }\mathrm{\Theta })\mathrm{\Delta }_Fx)(e^v,e^w)=(\mathrm{\Delta }_{_{loc}[G]}\mathrm{\Theta }(x))(e^v,e^w),$$ i.e. the identification is indeed as Hopf algebras to the relevant order. The right hand side here is $`\mathrm{\Theta }(x)(e^ve^w)=\mathrm{\Theta }(x)(e^{C(v,w)})=\mathrm{}^1x,C(v,w)`$ where $`C(v,w)`$ is the Campbell-Baker-Hausdorf series. Meanwhile, evaluating the coproduct $`\mathrm{\Delta }_F`$ we have $`((\mathrm{\Theta }\mathrm{\Theta })\mathrm{\Delta }_Fx)(e^v,e^w)=\mathrm{}^1(x,v+x,w)+\mathrm{}^1\alpha [x,e_i],ve^i,w+\mathrm{}^1\beta e^i,v[x,e_i],w`$ $`+\mathrm{}^1{\displaystyle \frac{\alpha ^2}{2}}v,[[x,e_i],e_j]e^i,we^j,w+\mathrm{}^1{\displaystyle \frac{\beta ^2}{2}}w,e^iw,e^jv,[[x,e_i],e_j]`$ $`+\mathrm{}^1({\displaystyle \frac{\alpha ^2}{2}}{\displaystyle \frac{1}{12}})v,e^iv,[x,e_j]w,[e^j,e_i]+\mathrm{}^1({\displaystyle \frac{\beta ^2}{2}}{\displaystyle \frac{1}{12}})v,e^i<v,e^jw,[[x,e_i],e_j]+\mathrm{}`$ $`=\mathrm{}^1x,v+w+[w,v](\alpha \beta )+{\displaystyle \frac{\alpha ^2}{2}}[w,[w,v]]+{\displaystyle \frac{\beta ^2}{2}}[v,[v,w]]+({\displaystyle \frac{\alpha ^2}{2}}{\displaystyle \frac{1}{12}})[[w,v],w]`$ $`+({\displaystyle \frac{\beta ^2}{2}}{\displaystyle \frac{1}{12}})[[v,w],v]+\mathrm{}`$ $`=\mathrm{}^1x,v+w+{\displaystyle \frac{1}{2}}[v,w]+{\displaystyle \frac{1}{12}}([v,[v,w]]+[[v,w],w])+\mathrm{}`$ for all $`x𝔤^{}`$. We used $`[x,v],w=xv,w=x,vw=x,[v,w]`$ and similarly for repeated commutators. Note that the increasing powers of $`\mathrm{}^1`$ with each evalution exactly match the increasing powers of $`\mathrm{}`$ in the powerseries coming from $`F`$. Hence the twisted coproduct reproduces the Campbell-Baker-Hausdorf series to low degree in its expansion. It is clear that this requirement and that we continue to reproduce the product of $`U_{\mathrm{}}(𝔤)`$ can be used to determine a universal formula for $`F,F^1`$, though again it would be beyond our scope to provide this here. ### 5.5. The Duflo map The Duflo map provides an independent check of our formula in Theorem 5.4.1 and gives some idea of the structure of $`F^1`$ at all orders. We recall \[Du\] that there is an invertible operator $`D`$ on $`S(𝔤)`$ defined by $$D=e^{_{k=1}^{\mathrm{}}\alpha _{2k}_{\mathrm{Tr}_{2k}}};\underset{k=1}{\overset{\mathrm{}}{}}\alpha _{2k}t^{2k}\frac{1}{2}\mathrm{ln}\left(\frac{\mathrm{sinh}(t/2)}{t/2}\right)$$ where $`_{\mathrm{Tr}_{2k}}`$ is the differential operator on $`S(𝔤)`$ given by the action of the element $`\mathrm{Tr}_{2k}=\mathrm{Tr}_𝔤((\mathrm{ad}_{})^{2k})S(𝔤^{})`$. The lowest order part is $$\alpha _2=\frac{1}{48},\alpha _4=\frac{1}{5760},_{\mathrm{Tr}_2}=\kappa _{ij}e^ie^j$$ in our conventions above. Duflo’s theorem is that when restricted to the ad-invariant subalgebra $`S(𝔤)^𝔤`$ the map $`\phi D`$ is an isomorphism of this with the centre $`Z(U(𝔤))`$. Now let $`F_{\mathrm{red}}^1U(𝔤^{})^2`$ be the effective $`F^1`$ when acting on invariants $`S(𝔤)^𝔤S(𝔤)^𝔤`$. This is given by normal ordering $`F^1`$ so that all terms have all elements of $`𝔤`$ to the right of all elements of $`𝔤^{}`$ (what we called $`::^R`$ above). Then project $`𝔤`$ to zero in the result because by definition it act by zero on invariant elements. The result is some power-series in $`U(𝔤^{})^2=S(𝔤^{})^2`$ since $`𝔤^{}`$ is being regarded as an Abelian Lie algebra. ###### Proposition 5.5.1. At least to $`O(\mathrm{}^3)`$, $`F_{\mathrm{red}}^1`$ for the cochain in Theorem 5.4.1 is a coboundary in the sense of \[M2, Chapter 2.3\] of the Duflo element, i.e., $$F_{\mathrm{red}}^1=(\mathrm{\Delta }\gamma )(\gamma ^1\gamma ^1),\gamma =e^{_{k=1}^{\mathrm{}}\alpha _{2k}\mathrm{}^{2k}\mathrm{Tr}_{2k}}$$ Here $`\gamma `$ viewed as an operator acting on $`S(𝔤)`$ is just $`D`$ in Duflo’s theorem after explicitly introducing the deformation scaling parameter. The coproduct $`\mathrm{\Delta }`$ is that of $`U(𝔤^{})`$. We expect this result to hold to all orders because $`F_{\mathrm{red}}^1`$ a coboundary of some cochain $`\gamma `$ implies that $$fg=\mu (F_{\mathrm{red}}^1.(fg))=\gamma .\mu (\gamma ^1.f\gamma ^1.g))=D(\mu (D^1fD^1g))$$ for all $`f,gS(𝔤)^𝔤`$, where $`D`$ denotes $`\gamma `$ acting on $`S(𝔤)`$. We used in the first equality that $`U(𝔤^{})`$ acts covariant on $`S(𝔤)`$ with its initial product $`\mu `$ and hence we can move the action of $`\mathrm{\Delta }\gamma `$ to the left as the action of $`\gamma `$. This means that the modified product restricted to invariant elements is an isomorphism of algebras (this is the meaning of $`F_{\mathrm{red}}^1`$ being a coboundary as explained in \[M2\]). In the light of Duflo’s theorem we explect $`F_{\mathrm{red}}^1`$ therefore to be a coboundary of an invertible element $`\gamma `$ whose action is the same as the operator $`D`$ in Duflo’s theorem. This leads to the statement of the proposition. We now verify the proposition to the order $`O(\mathrm{}^3)`$ available to us. In the expression in Theorem 5.4.1 all the terms have some $`e_i`$ already to the right and therefore fail to contribute, except $`:Q_1Q_2:`$ and the $`\kappa _{ij}e^ie^j`$ terms. We write the former using $$e_je^i=e^ie_j+(e^ie_j)=e^ie_j+f_{jki}e^k$$ where $`[e_i,e_j]=f_{ijk}e_k`$ defines the structure constants. Then $$:Q_1Q_2:=e_ie^je_je^if_{jki}f_{imj}e^ke^mf_{kji}f_{mij}e^ke^m\kappa _{ij}e^ie^j$$ discarding terms acting trivially on invariant elements. As a result we have $$F_{\mathrm{red}}^1=11\frac{\mathrm{}^2}{24}\kappa _{ij}e^ie^j+O(\mathrm{}^3)$$ (in fact the next term should be $`O(\mathrm{}^4)`$). On the other hand from the above $$\gamma =e^{\frac{\mathrm{}^2}{48}c+O(\mathrm{}^4)};c=\kappa _{ij}e^ie^j,(\mathrm{\Delta }\gamma )(\gamma ^1\gamma ^1)=e^{\frac{\mathrm{}^2}{48}(\mathrm{\Delta }cc11c)+O(\mathrm{}^4)}=F_{\mathrm{red}}^1$$ to lowest order. The same proposition would provide a check to all orders of any cochain found. At the moment we have provided a check of our order $`O(\mathrm{}^3)`$ result. ### 5.6. Example: noncommutative Minkowski space as cochain twist. Here we verify (33) for the algebra $`[t,x_i]=\mathrm{}x_i`$ which has been proposed as noncommutative spacetime (the so-called bicrossproduct model). For convenience we take only one $`x=x_i`$ rather than $`i=1,2,3`$ for spacetime, however the structure is exactly similar. In this case we exhibit a candidate for $`F^1`$ to the next order, i.e. up to $`O(\mathrm{}^4)`$. In this model there is a representation of the algebra in terms of $`2\times 2`$ matrices, as $`t\left(\begin{array}{cc}\mathrm{}/2& 0\\ 0& \mathrm{}/2\end{array}\right),x\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right).`$ These matrices can be exponentiated to give $`\mathrm{exp}(pt+qx)`$ $`=`$ $`\left(\begin{array}{cc}e^{\frac{\mathrm{}p}{2}}& \frac{\left(1+e^\mathrm{}p\right)q}{e^{\frac{\mathrm{}p}{2}}\mathrm{}p}\\ 0& e^{\frac{\left(\mathrm{}p\right)}{2}}\end{array}\right).`$ A little matrix multiplication shows that $`\mathrm{exp}(pt+qx).\mathrm{exp}(rt+sx)`$ $`=`$ $`\mathrm{exp}\left((p+r)t+{\displaystyle \frac{\left(p+r\right)\left(\left(1+e^\mathrm{}p\right)qr+e^\mathrm{}p\left(1+e^\mathrm{}r\right)ps\right)}{\left(1+e^{\mathrm{}\left(p+r\right)}\right)pr}}x\right),`$ so in this case we have a closed form for the CBH formula. Further calculation with this algebra gives $`\phi (x^nt^m)`$ $`=`$ $`x^n(t^m+{\displaystyle \frac{nm}{2}}\mathrm{}t^{m1}+{\displaystyle \frac{n(3n+1)m(m1)}{24}}\mathrm{}^2t^{m2}`$ $`+{\displaystyle \frac{n^2(n+1)m(m1)(m2)}{48}}\mathrm{}^3t^{m3}+O(\mathrm{}^4)).`$ From this we can calculate $`\phi ^1(x^nt^m)`$ $`=`$ $`x^n(t^m{\displaystyle \frac{nm}{2}}\mathrm{}t^{m1}+{\displaystyle \frac{n(3n1)m(m1)}{24}}\mathrm{}^2t^{m2}`$ $`+{\displaystyle \frac{n^2(1n)m(m1)(m2)}{48}}\mathrm{}^3t^{m3}+O(\mathrm{}^4)).`$ If we combine this with the following formula for multiplication in $`U_{\mathrm{}}`$, $`t^sx^r={\displaystyle \underset{p=0}{\overset{s}{}}}C_p^s(r\mathrm{})^px^rt^{sp},`$ (where $`C_p^s`$ is a binomial coefficient) we get the formula (44) $`(x^nt^m)(x^rt^s)`$ $`=`$ $`x^{n+r}t^{m+s}+{\displaystyle \frac{\mathrm{}}{2}}(mrns)x^{n+r}t^{m+s1}`$ $`+{\displaystyle \frac{\mathrm{}^2}{24}}(mrm^2r3mr^2+3m^2r^2+ns2mns3n^2s2mrs6mnrs`$ $`ns^2+3n^2s^2)x^{n+r}t^{m+s2}+{\displaystyle \frac{\mathrm{}^3}{48}}(2mr^2+3m^2r^2m^3r^2+2mr^3`$ $`3m^2r^3+m^3r^3+2n^2s2mn^2s2n^3sm^2nrs3mn^2rs+2mr^2s`$ $`2m^2r^2s+3mnr^2s3m^2nr^2s3n^2s^2+2mn^2s^2+3n^3s^2+mnrs^2`$ $`+3mn^2rs^2+n^2s^3n^3s^3)x^{n+r}t^{m+s3}+O(\mathrm{}^4).`$ On $`S(𝔤)`$, $`\mathrm{ad}_t`$ is identified with $`x\frac{\mathrm{d}}{\mathrm{d}x}`$, and $`\mathrm{ad}_x`$ is identified with $`x\frac{\mathrm{d}}{\mathrm{d}t}`$. We take the dual basis $`\widehat{t},\widehat{x}𝔤^{}`$. Then on $`S(𝔤)`$, $`\widehat{t}`$ is identified with $`\frac{\mathrm{d}}{\mathrm{d}t}`$, and $`\widehat{x}`$ is identified with $`\frac{\mathrm{d}}{\mathrm{d}x}`$. Using this, it can be explicitly checked that (33) gives the deformed multiplication for this algebra up to $`O(\mathrm{}^3)`$. The third order part of (44) can be given by $`G^{(3)}=(e_ie^je^ke_ke_je^ie_ke^je_ie^ie_je^k2e^ie_je^ke_ie^je_k)/96,`$ where we sum over $`i,j,k`$. Note, however, that this expression is not unique in the same manner that (33) at order $`\mathrm{}^2`$ is not unique as we have seen. With more work one may exploit the non-uniqueness and expect to achieve the features in Section 5.4 with respect to the coproduct $`\mathrm{\Delta }_F`$ as well. Note that we do not necessarily expect a unique $`F,F^1`$ for any given Lie algebra (the uniqueness proposed in Section 5.4 was for a universal $`F,F^1`$ applicable to all Lie algebras). ## 6. Mackey quantisations $`C^{\mathrm{}}(N)>U_{\mathrm{}}(g)`$ of Homogeneous spaces as cochain twists Suppose that a Lie group $`G`$ with Lie algabra $`𝔤`$ acts on a manifold $`N`$. In this case there is a standard ’quantisation’ for the system due to Mackey and much used in physics, in which the initial algebra is $`C^{\mathrm{}}(N)S(𝔤)C^{\mathrm{}}(N\times 𝔤^{})`$ (i.e. functions polynomial in the $`𝔤^{}`$ direction). This is deformed or quantised to the cross product $`C^{\mathrm{}}(N)>U_{\mathrm{}}(𝔤)`$. Here, for $`v𝔤`$ and $`fC^{\mathrm{}}(N)`$, $`(v_{\mathrm{}}f)(x)=\mathrm{}f^{}(x;v(x))`$, and $`U_{\mathrm{}}(𝔤)`$ has the relation $`vwwv=\mathrm{}[v,w]`$ in terms of the Lie bracket $`[,]`$ on $`𝔤`$. This algebra acts on the $`L^2`$ sections of a bundle whose fiber over $`xN`$ is a representation of the stabiliser of $`x`$ in $`G`$. In this section we show that the results of Section 5 may be extended to this case also. The theory here reduces to that of Section 5 when $`N`$ is a point. Note that $`M=N\times 𝔤^{}`$ is indeed a Poisson manifold, because the quantisation above can be viewed as a flat deformation. Its Poisson bracket has a semidirect product form $$\{f,g\}=0,\{v,f\}=vf,\{v,w\}=[v,w]$$ for $`f,gC^{\mathrm{}}(N)`$ and $`v,w𝔤^{}`$. Our goal is to lift this Poisson bivector to an element of a suitable $``$ and hence to a cochain $`F`$ at least to order $`O(\mathrm{}^2)`$, i.e. to express the Mackey quantisation as a cochain twist. ### 6.1. To first order We will be extending the results from the previous ‘CBH’ case in Section 5; we denote the cochain components there by $`G_{CBH}^{(i)}`$. ###### Definition 6.1.1. Take a dual basis $`(e_i,e^i)`$ with $`e_i𝔤`$ and $`e^i𝔤^{}`$. Then define some vector fields on $`M=N\times 𝔤^{}`$ by the following formulae, where $`v𝔤S(𝔤)`$ and $`gC^{\mathrm{}}(N)`$. $`\stackrel{ˇ}{e}_i(v)=[e_i,v]`$ , $`\stackrel{ˇ}{e}_i(g)=\mathrm{\hspace{0.17em}0},`$ $`\stackrel{ˇ}{e}^i(v)=e^i(v)`$ , $`\stackrel{ˇ}{e}^i(g)=\mathrm{\hspace{0.17em}0},`$ $`\stackrel{ˇ}{c}_i(v)=\mathrm{\hspace{0.17em}0}`$ , $`\stackrel{ˇ}{c}_i(g)=e_ig.`$ ###### Proposition 6.1.2. The Poisson structure described is given by $`G^{(1)}=(\stackrel{ˇ}{c}_i+\stackrel{ˇ}{e}_i/2)\stackrel{ˇ}{e}^i.`$ Proof: Recall that for $`G^{(1)}=XY`$ we have $`\{a,b\}=X(a)Y(b)Y(a)X(b)`$. Hence for $`v,w𝔤`$ and $`g,kC^{\mathrm{}}(N)`$: $`\stackrel{ˇ}{e}_i(v)\stackrel{ˇ}{e}^i(w)\stackrel{ˇ}{e}^i(v)\stackrel{ˇ}{e}_i(w)`$ $`=`$ $`[e_i,v]e^i(w)e^i(v)[e_i,w]=[w,v][v,w]=\mathrm{\hspace{0.17em}2}[w,v],`$ $`\stackrel{ˇ}{e}_i(v)\stackrel{ˇ}{e}^i(g)\stackrel{ˇ}{e}^i(v)\stackrel{ˇ}{e}_i(g)`$ $`=`$ $`0,`$ $`\stackrel{ˇ}{e}_i(g)\stackrel{ˇ}{e}^i(k)\stackrel{ˇ}{e}^i(g)\stackrel{ˇ}{e}_i(k)`$ $`=`$ $`0,`$ $`\stackrel{ˇ}{c}_i(v)\stackrel{ˇ}{e}^i(w)\stackrel{ˇ}{e}^i(v)\stackrel{ˇ}{c}_i(w)`$ $`=`$ $`0,`$ $`\stackrel{ˇ}{c}_i(v)\stackrel{ˇ}{e}^i(g)\stackrel{ˇ}{e}^i(v)\stackrel{ˇ}{c}_i(g)`$ $`=`$ $`e^i(v)(e_ig)=vg,`$ $`\stackrel{ˇ}{c}_i(g)\stackrel{ˇ}{e}^i(k)\stackrel{ˇ}{e}^i(g)\stackrel{ˇ}{c}_i(k)`$ $`=`$ $`0`$ as desired. $`\mathrm{}`$ ###### Proposition 6.1.3. The Lie brackets between the $`\stackrel{ˇ}{e}_i`$, $`\stackrel{ˇ}{c}_i`$ and $`\stackrel{ˇ}{e}^i`$ are given as follows. The $`\stackrel{ˇ}{c}_i`$ commute with both the $`\stackrel{ˇ}{e}_j`$ and the $`\stackrel{ˇ}{e}^j`$. The $`\stackrel{ˇ}{e}_i`$ have the usual Lie bracket for $`𝔤`$ among themselves. The $`\stackrel{ˇ}{e}^i`$ commute among themselves. The $`\stackrel{ˇ}{c}_i`$ have the usual Lie bracket for $`𝔤`$ among themselves. The bracket of the $`e_i`$ with the $`e^j`$ is given by the coadjoint action. Thus the given fields form a Lie algebra $`=𝔤<𝔤^{}𝔤`$. Proof: Check against $`v`$ and $`g`$. The most difficult ones are: $`[\stackrel{ˇ}{c}_i,\stackrel{ˇ}{c}_j](g)`$ $`=`$ $`(e_i(e_jg)e_i(e_jg))=([e_i,e_j]g),`$ $`[\stackrel{ˇ}{e}_i,\stackrel{ˇ}{e}^j](v)`$ $`=`$ $`\stackrel{ˇ}{e}^j([e_i,v])=e^j([e_i,v]).`$ The last equation shows the coadjoint action, $`\mathrm{coad}_w(\psi )=\psi \mathrm{ad}_w`$. $`\mathrm{}`$ We shall use this Lie algebra to induce the Mackey quantisation. We have already seen above that this is sufficient at order $`\mathrm{}`$. ###### Proposition 6.1.4. The preconnection is for $`v,w𝔤`$ and $`f,gC^{\mathrm{}}(N)`$: $`\widehat{}_v\mathrm{d}w`$ $`=`$ $`\mathrm{d}[v,w]/2,`$ $`\widehat{}_v\mathrm{d}g`$ $`=`$ $`\mathrm{d}(vg),`$ $`\widehat{}_f\mathrm{d}w`$ $`=`$ $`0,`$ $`\widehat{}_f\mathrm{d}g`$ $`=`$ $`0.`$ Proof: We use the formula involving Lie derivatives; $`\widehat{}_a\xi `$ $`=`$ $`\stackrel{ˇ}{e}^i(a)_{\stackrel{ˇ}{c}_i+\stackrel{ˇ}{e}_i/2}\xi (\stackrel{ˇ}{c}_i+\stackrel{ˇ}{e}_i/2)(a)_{\stackrel{ˇ}{e}^i}\xi .`$ This gives $`\widehat{}_v\mathrm{d}w`$ $`=`$ $`e^i(v)\mathrm{d}(\stackrel{ˇ}{c}_i+\stackrel{ˇ}{e}_i/2)(w)[e_i,v]\mathrm{d}(\stackrel{ˇ}{e}^i)(w)/2`$ $`=`$ $`e^i(v)\mathrm{d}[e_i,w]/2,`$ $`\widehat{}_v\mathrm{d}g`$ $`=`$ $`e^i(v)\mathrm{d}(\stackrel{ˇ}{c}_i+\stackrel{ˇ}{e}_i/2)(g)[e_i,v]\mathrm{d}(\stackrel{ˇ}{e}^i)(g)/2`$ $`=`$ $`e^i(v)\mathrm{d}(e_ig),`$ $`\widehat{}_f\mathrm{d}w`$ $`=`$ $`(e_if)\mathrm{d}(\stackrel{ˇ}{e}_i)(w)=\mathrm{\hspace{0.17em}0},`$ $`\widehat{}_f\mathrm{d}g`$ $`=`$ $`(e_if)\mathrm{d}(\stackrel{ˇ}{e}_i)(g)=\mathrm{\hspace{0.17em}0}`$ as required. $`\mathrm{}`$ ###### Proposition 6.1.5. The curvature is given by, for $`v,w,z𝔤`$ and $`f,g,hC^{\mathrm{}}(N)`$: $`R(v,w)(\mathrm{d}z)`$ $`=`$ $`\mathrm{d}([[v,w],z])/4,`$ $`R(v,w)(\mathrm{d}g)`$ $`=`$ $`0,`$ $`R(v,g)(\mathrm{d}z)`$ $`=`$ $`R(v,g)(\mathrm{d}h)=\mathrm{\hspace{0.17em}0},`$ $`R(f,g)(\mathrm{d}z)`$ $`=`$ $`R(f,g)(\mathrm{d}h)=\mathrm{\hspace{0.17em}0}.`$ Proof: First we need to state the Lie brackets of the vector fields, using $`[\widehat{a},\widehat{b}]=\widehat{\{a,b\}}`$: $`[\widehat{v},\widehat{w}]`$ $`=`$ $`\widehat{\{v,w\}}=\widehat{[v,w]},`$ $`[\widehat{v},\widehat{g}]`$ $`=`$ $`\widehat{\{v,g\}}=\widehat{(vg)},`$ $`[\widehat{f},\widehat{g}]`$ $`=`$ $`\widehat{\{f,g\}}=\mathrm{\hspace{0.17em}0}.`$ Now we find the curvatures. For $`R(v,w)(\mathrm{d}z)`$ the computation is as in Proposition 5.1.1 with the same result. For the new case we have $$R(v,w)(\mathrm{d}g)=\mathrm{d}(v(wg)w(vg)[v,w]g)=0$$ since $``$ is a representation of the Lie algebra (if it were a cocycle representation we would have curvature). In the expressions for $`R(v,g)`$ and $`R(f,g)`$ every term is individually zero.$`\mathrm{}`$ ### 6.2. To second order The multiplication on $`C^{\mathrm{}}(N)>U_{\mathrm{}}(𝔤)`$ is given by $`(f\underset{¯}{v})(g\underset{¯}{w})`$ $`=`$ $`f.(\underset{¯}{v}_{(1)}_{\mathrm{}}g)\underset{¯}{v}_{(2)}\underset{¯}{w},`$ where the second product is in $`U_{\mathrm{}}(g)`$ and the coproduct is the usual $`\mathrm{\Delta }v=v1+1v`$ for $`v𝔤`$. In terms of deformations of $`C^{\mathrm{}}(N)S(𝔤)`$ we can decompose the order $`\mathrm{}^2`$ part of the product as (45) $`f.g(\mathrm{}^2\mathrm{part}\mathrm{of}\underset{¯}{v}\underset{¯}{w})+f.(\mathrm{}\mathrm{part}\mathrm{of}\underset{¯}{v}_{(1)}_{\mathrm{}}g)(\mathrm{}\mathrm{part}\mathrm{of}\underset{¯}{v}_{(2)}\underset{¯}{w})`$ (46) $`+f.(\mathrm{}^2\mathrm{part}\mathrm{of}\underset{¯}{v}_{(1)}_{\mathrm{}}g)(\mathrm{}^0\mathrm{part}\mathrm{of}\underset{¯}{v}_{(2)}\underset{¯}{w}).`$ The first term of (45) is given by $`\mathrm{}^2(1G_{CBH}^{(2)}{}_{1}{}^{})(1G_{CBH}^{(2)}{}_{2}{}^{})`$, where the final suffices $`1,2`$ denote the two pieces of $`G_{CBH}^{(2)}`$ (summation understood). The second term of (45) is, where hat denotes ommission, $`\mathrm{}f(x).{\displaystyle \underset{i}{}}(v_ig)(x)(\mathrm{}\mathrm{part}\mathrm{of}(v_1\mathrm{}\widehat{v}_i\mathrm{}v_n)(w_1\mathrm{}w_m)).`$ We can separate this into two stages, first the moving the $`v_i`$ stage, and then the $``$ multiplication. The first is given by $`\mathrm{}(1\stackrel{ˇ}{e}^i)(\stackrel{ˇ}{e}_i1)`$, and the second by $`\mathrm{}(1G_{CBH}^{(2)}{}_{1}{}^{})(1G_{CBH}^{(2)}{}_{2}{}^{})`$ as above. The third term of (45) is, $`\mathrm{}^2f.{\displaystyle \underset{i<j}{}}(v_iv_jg)v_1\mathrm{}\widehat{v}_i\mathrm{}\widehat{v}_j\mathrm{}v_nw_1\mathrm{}w_m.`$ This is given by $`\frac{1}{2}\mathrm{}^2(1\stackrel{ˇ}{e}^i\stackrel{ˇ}{e}^j)(1\stackrel{ˇ}{e}_i\stackrel{ˇ}{e}_j)`$, giving in total $`G^{(2)}=(1G_{CBH}^{(2)}{}_{1}{}^{})(1G_{CBH}^{(2)}{}_{2}{}^{})+{\displaystyle \underset{i}{}}(1G_{CBH}^{(1)}{}_{1}{}^{}\stackrel{ˇ}{e}_{}^{i})(\stackrel{ˇ}{e}_iG_{CBH}^{(1)}{}_{2}{}^{})+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}(1\stackrel{ˇ}{e}^i\stackrel{ˇ}{e}^j)(1\stackrel{ˇ}{e}_i\stackrel{ˇ}{e}_j)`$ A special case is of course $`N=G`$ and action by left translation. Then the Mackey quantisation $`C^{\mathrm{}}(G)>U_{\mathrm{}}(𝔤)`$ is a quantisation of $`T^{}G=G\times 𝔤^{}`$ and the Poisson-bracket above becomes the standard sympletic structure on $`T^{}G`$. In that case we have an actual connection $``$ in Section 6.1. ### 6.3. Special case of $`T^{}G`$ In general we have a Poisson map $`T^{}NN\times 𝔤^{}`$ defined using the moment map by $`(n,p)(n,x_{}(n),p)`$ where $`x_\xi `$ is the vector field for the action of $`\xi 𝔤`$ on $`N`$. This means a map $$C^{\mathrm{}}(N\times 𝔤^{})C^{\mathrm{}}(T^{}N)$$ which will be surjective in the case that the action is locally transitive. In this way the Mackey quantisation results above can in principle induce quantisations of $`T^{}N`$. In terms of functions on $`T^{}N`$, $`fC^{\mathrm{}}(N)`$ corresponds to $`\overline{f}=\pi ^{}fC^{\mathrm{}}(T^{}N)`$ and $`v𝔤`$ corresponds to the function $`\overline{v}(x,p)=p,v(x)`$. The elements of the algebra we order putting all elements of $`𝔤`$ to the right. We get the relation $`\overline{v}\overline{f}=\mathrm{}\overline{vf}+\overline{f}\overline{v}`$. In terms of commutators, $`[\overline{v},\overline{f}]=\mathrm{}\overline{vf}`$, and we would like this to be given by a Poisson bracket on $`T^{}N`$. This means $`\omega (\mathrm{d}\overline{v},\mathrm{d}\overline{f})=\overline{vf}`$, or in $`(x,p)`$ coordinates $`{\displaystyle \frac{\overline{f}}{x^j}}v(x)^j`$ $`=`$ $`\omega ^{(i+n)j}{\displaystyle \frac{\overline{v}}{p_i}}{\displaystyle \frac{\overline{f}}{x^j}}+\omega ^{ij}{\displaystyle \frac{\overline{v}}{x^i}}{\displaystyle \frac{\overline{f}}{x^j}}`$ $`=`$ $`\omega ^{(i+n)j}v(x)^i{\displaystyle \frac{\overline{f}}{x^j}}+\omega ^{ij}p_k{\displaystyle \frac{v^k}{x^i}}{\displaystyle \frac{\overline{f}}{x^j}}.`$ Provided the vector fields $`v(x)`$ for the action of $`𝔤`$ span the tangent space at each point, this implies that $`\omega ^{(j+n)j}=1=\omega ^{j(j+n)}`$ and all others are zero, i.e. the standard symplectic form on $`T^{}N`$. Now we calculate $`\omega (\mathrm{d}\overline{v},\mathrm{d}\overline{w})`$ $`=`$ $`v^jp_k{\displaystyle \frac{w^k}{x^j}}w^jp_k{\displaystyle \frac{v^k}{x^j}}=\overline{[v,w]}.`$ These conditions are met for $`T^{}G`$ with Mackey quantisation $`C^{\mathrm{}}(G)>U_{\mathrm{}}(𝔤)`$ where the action is by left translation. Hence in this case the formulae in Section 6.1 define an actual compatible connection $``$ on $`T^{}G`$. This model will be investigated further elsewhere. ## 7. Quantum groups $`_q[G]`$ and related examples Here, for completeness, we show that Drinfelds original construction of quantum groups $`_q[G]`$ (this means more precisely the dual of Drinfeld’s construction) can also be formulated as a cochain module twist. This is not fundamentally new but a useful point of view that motivated the above. Actually, this is a general observation for any twist, in the setting of Section 2. Thus, if we are are interested in $`A`$ an initial undeformed Hopf algebra, let $`A^{}`$ be a dually paired Hopf algebra and $`H=A^{}A^{}^{\mathrm{op}}`$. The use of a dual here is of course avoided if we work with comodule twists rather than module ones. In our case $`A^{}`$ acts on $`A`$ from the left by $`ha=(\mathrm{id}h)(\mathrm{\Delta }a)`$ and $`A^{}^{\mathrm{op}}`$ acts by $`ha=(h\mathrm{id})(\mathrm{\Delta }a)`$. In this way $`A`$ becomes $`H`$-covariant as an algebra. Now let $`fA^{}A^{}`$ be a cochain. This induces a cochain $$F=f_{13}f_{24}^1HH$$ where the suffices refer to the position in the four-fold tensor power of $`A^{}`$. The modified algebra $`A_F`$ induced by this cochain has product $$ab=(f_{24}f_{13}^1(ab))=f(a_{(1)},b_{(1)})a_{(2)}b_{(2)}f^1(a_{(3)},b_{(3)})$$ which from the Drinfeld point of view is a usual (co)twist of the Hopf algebra $`A`$ into a coquasiHopf algebra. Moreover, $`A_F`$ will be covariant under $`H_F`$ with coproduct $$\mathrm{\Delta }_F(hg)=f_{13}f_{24}^1\mathrm{\Delta }_{13}(h)\mathrm{\Delta }_{24}(g)f_{24}f_{13}^1$$ where the products are in the square of $`A^{}A^{}^{\mathrm{op}}`$. If we denote by $`A_f^{}`$ the usual twist of $`A^{}`$ by $`f`$ with coproduct $`\mathrm{\Delta }_f(h)=f(\mathrm{\Delta }h)f^1`$. We see that $`H_F=A_f^{}(A_f^{})^{\mathrm{op}}`$, all as potentially quasi-Hopf algebras. This is the general situation and amounts to a recasting of the standard Drinfeld (co)twist of $`A`$ as a module-algebra cochain twist for a suitably doubled up $`H`$. We can of course apply this to Drinfelds example where $`A=[G]`$ is a suitable form of the coordinate ring of a classical simple Lie group with Lie algebra $`𝔤`$, and $`A^{}=U(𝔤)`$. It is understood that we extend the above to allow formal powerseries in a parameter $`\mathrm{}`$ (with $`q=e^\frac{\mathrm{}}{2}`$). In this case Drinfeld showed the existence (as a formal power-series) of a suitable $`f`$ such that $`A_f^{}U_q(𝔤)`$ as a Hopf algebra. Here the coboundary $`\varphi _f`$ and hence the associator $`\mathrm{\Phi }_f`$ are nontrivial but $`\varphi _f`$ is central in the sense that $`A_f^{}`$ remains an ordinary Hopf algebra. This translates in the above reformulation into the statement that $`F`$ and its coboundary $`\varphi _FHHH`$ are nontrivial but that $`\varphi _F`$ acts trivially on $`A`$, so that $`A_F_q[G]`$ remains associative, namely the usual quantum group coordinate algebra. We have $`=𝔤𝔤^{\mathrm{op}}`$ as the incuding Lie algebra from our current point of view. We can further construct $`\mathrm{\Omega }(_q[G])=\mathrm{\Omega }(G)_F`$ and compute its associated preconnection and curvature (which is nonzero). These results are the same as those presented in \[BM1\] (albeit from a different supercoquasiHopf algebra point of view) so we do not repeat them here. Note that Drinfeld’s $`f`$ is not known very explicitly (except at lowest order where the Poisson-bracket induced by the above is the usual Drinfeld-Sklyanin one), however its existence holds very generally and in a canonical way for quantum-group related examples. Moreover, it can happen that $`H=U_q(𝔤)`$ and $`F=f`$ may again have $`\varphi _F`$ acting trivially on a particular classical algebra, such as on a highest weight orbit. Here $`=𝔤`$ is the inducing Lie algebra and we use Drinfeld’s cochain without any doubling. This was the case in \[DGM\] where it was an associative quantum sphere was constructed in this way with $`=su_2`$. This is therefore an early example of the cochain-quantisation method genuinely used. ## 8. Hidden nonassociativity In the above we have recovered, at least to some order, several standard associative quantum algebras of interest in physics as cochain twists (we do not just mean q-deformed or quantum group examples). Here the cochains are not required to be cocycles and this relaxation appears to be necessary. It means, however, that even though the algebra of ’functions’ happens to remain associative, there is an underlying nonassociativity behind the scenes in all these quantum algebras. We now turn to this aspect. First of all, as our algebras become quantized, their covariance Lie algebra $``$ gets deformed to a quasi-quantum group $`U()_F`$ as explained in Section 2. These are in principle ’noncoassociative’ and are looked at for our various examples in Section 8.1. Next, our quantum algebras are all equivalent in a certain monoidal categorical sense to the unquantised algebras, with the result that not only the algebras but all functorial constructions on them are similarly quantised, for example differential forms on the classical phase space and the Dirac operator deform naturally to the quantum algebras, but nonassociatively. We consider these in Sections 8.2 and 8.3 respectively. ### 8.1. The quasiHopf algebras $`U()_F`$ From Section 2, the deformed algebra $`A_F`$ remains covariant, but under the quasi-Hopf algebra $`H_F`$. In our cases of interest $`H=U()`$ which is also the algebra of $`H_F`$. Its coproduct, however, is modified to $$\mathrm{\Delta }_F(X)=F(X1+1X)F^1=X1+1X\mathrm{}[G^{(1)},X1+1X]+O(\mathrm{}^2)$$ for any $`X`$. The leading order data here defines a quasi-Lie bialgebra $`(,\delta ,\psi )`$ where $$\delta X=[X1+1X,G^{(1)}]=\mathrm{ad}_X(G^{(1)}).$$ If it happens that $`\delta `$ obeys the cojacobi identity $$(\delta \mathrm{id})\delta X+\mathrm{cyclic}=0$$ then we have an ordinary Lie bialgebra, which means that at lowest order at least, $`U()_F`$ remains an ordinary (not quasi) Hopf algebra. This is already the case for the Drinfeld twist examples in Section 7 and indeed the $`O(\mathrm{}^2)`$ part $`\psi `$ of $`\varphi `$ is a multiple of the ad-invariant Cartan tensor $`n\mathrm{\Lambda }^3(𝔤)`$ defined by the Killing form. In general the cojacobiator above is given by $`\mathrm{ad}_X(\psi )`$ and this is the fundamental reason why the covariance algebra remains (co)associative for such examples. But let us see how the situation fares for our non-quantum group examples. Thus, $`\psi `$ is given for the $`=so(1,3)`$ example in Section 4.5 and from the expression there one may readily compute that $$\mathrm{ad}_{X_i}(\psi )0,\mathrm{ad}_{Y_i}(\psi )=0$$ i.e. rotationally invariant (as to be expected as the whole construction is) but not invariant under boosts. Thus our ‘sphere at infinity’ example in Section 4 gives us a quasi-Hopf algebra version of $`=so(1,3)`$. Next up, we CBH or $`U_{\mathrm{}}(g)`$ example in Section 5, we have seen in Section 5.4 what the twisted coproduct $`\mathrm{\Delta }_F`$ looks like to $`O(\mathrm{}^3)`$ when acting on $`[𝔤]H=U(𝔤<𝔤^{})=U(𝔤)<[𝔤]`$. We have seen that this twists to a local form of the classical coordinate ring $`[G]`$. On the other hand the coproduct of $`U(𝔤)`$ remains unchanged after twisting to this order because all elements are $`g`$-invariant under commutator in the bigger algebra and hence $$[\mathrm{\Delta }v,G^{(i)}]=[v1+1v,G^{(i)}]=[v,G_1^{(i)}]G_2^{(i)}+G_1^{(i)}[v,G_2^{(i)}]=0.$$ This is clear since the expressions involve only paired bases and dual basis, the Killing form etc. and the commutators of $`v𝔤`$ are given by the adjoint and coadjoint actions. Such invariance would be a reasonable requirement to all orders for any universal formula for $`F,F^1`$. Since $`H=U(𝔤)<[𝔤]`$ is generated by $`U(𝔤)`$, $`[𝔤]`$ and has the same algebra after twisting, we conclude that $$U(g<g^{})_F=U(g)<[g]_F$$ as an ordinary Hopf algebra. Moreover, this is locally isomorphic to $`U(g)<[G]=D(U(g))`$, the Drinfeld quantum double, which is an ordinary Hopf algebra. It is known that $`U_{\mathrm{}}(𝔤)`$ is always covariant under $`D(U(𝔤))`$ and this was explored for $`U_{\mathrm{}}(su_2)`$ (the so-called universal ‘fuzzy sphere’) in \[BaMa\]. In this case the background covariance becomes a quantum group covariance but remains associative (there is still hidden nonassociativity, see below). Finally, the Mackey case is a nontrivial extension of the CBH case and has a larger symmetry group. We do not make the full analysis here but suffice it to say that one reason that the CBH case works from an algebraic point of view to the fact that $`U_{\mathrm{}}(𝔤)`$ is a (cocommutative) Hopf algebra and therefore has a larger quantum group covariance based on the Drinfeld double; there is no such argument for the general Mackey case but there are cases when $`N`$ is itself a group and acts back on $`G`$ such that the Mackey quantisation in an algebraic form becomes a bicrossproduct Hopf algebra, see \[M2\]. In such cases one might expect similar behaviour to the CBH case above, but not in general. The bicrossproduct case includes the deformed Poincaré group for the noncommutative spacetimes mentioned in Section 5.6, see \[MR\]. ### 8.2. Quasiassociative quantum differential calculi Next, by applying the same cochain twist to the classical exterior algebra, we obtain noncommutative differential calculi on our various quantisations. We mean differential calculus in the sense of noncommutative geometry but in a monoidal weakly associative category, and we will see that our calculi on these examples are indeed nonassociative. They do, however, have the merit of classical dimensions in each degree. Again, the quantum group case in Section 7 was already covered in an equivalent form in \[BM1\] with the main result that the resulting $`\mathrm{\Omega }(_q[G])`$ have curvature and hence are not associative even though the quantised algebra happens to be. But let us see how our non-quantum group-related examples fare. First, for our sphere at infinity example. The given vector fields act by Lie derivative, which commutes with the $`\mathrm{d}`$ operator. Some large calculations give the special cases: $`x\mathrm{d}x`$ $`=`$ $`x.\mathrm{d}x+{\displaystyle \frac{\mathrm{}}{2z}}(x^2y.\mathrm{d}x+x(1x^2).\mathrm{d}y)+{\displaystyle \frac{\mathrm{}^2}{8}}x.\mathrm{d}x+O(\mathrm{}^3),`$ $`x\mathrm{d}y`$ $`=`$ $`x.\mathrm{d}y+{\displaystyle \frac{\mathrm{}}{2z}}(xy^2.\mathrm{d}x+y(1x^2).\mathrm{d}y)+{\displaystyle \frac{\mathrm{}^2}{8}}x.\mathrm{d}y+O(\mathrm{}^3),`$ $`y\mathrm{d}x`$ $`=`$ $`y.\mathrm{d}x{\displaystyle \frac{\mathrm{}}{2z}}(x(1y^2).\mathrm{d}x+x^2y.\mathrm{d}y)+{\displaystyle \frac{\mathrm{}^2}{8}}y.\mathrm{d}x+O(\mathrm{}^3),`$ $`y\mathrm{d}y`$ $`=`$ $`y.\mathrm{d}y{\displaystyle \frac{\mathrm{}}{2z}}(y(1y^2).\mathrm{d}x+xy^2.\mathrm{d}y)+{\displaystyle \frac{\mathrm{}^2}{8}}y.\mathrm{d}y+O(\mathrm{}^3).`$ Moreover, since the connection $``$ arising from the noncommutativity of the calculus at lowest degree turned out to be the Levi-Civita one, and since this has (constant) curvature, we know that the exterior algebra of this quantised sphere is necessarily nonassociative. Next, for the CBH or $`U_{\mathrm{}}(g)`$ example, again we have found expressions for the curvature in terms of a double commutators. Whether or not this vanishes depends on the Lie algebra in question: for the Heisenberg Lie algebra for example, one has $`R=0`$. However, for a simple Lie algebra such double commutators will not vanish and there is curvature, hence nonassociativity of $`\mathrm{\Omega }(U_{\mathrm{}}(𝔤))`$. One can also write the deformed calculus explicitly: $$v\mathrm{d}w=v\mathrm{d}w+\mathrm{}\beta (e^iv)\mathrm{d}e_iw)+O(\mathrm{}^3)=v\mathrm{d}w+\mathrm{}\beta \mathrm{d}[v,w]+O(\mathrm{}^3)$$ $$\mathrm{d}wv=(\mathrm{d}w)v+\mathrm{}\alpha (\mathrm{d}e_iw)e^iv+O(\mathrm{}^3)=(\mathrm{d}w)v+\mathrm{}\alpha \mathrm{d}[v,w]+O(\mathrm{}^3)$$ where we use the expression for $`F^1`$ in Section 5.4. The second order terms fail to contribute because $`e^iv=v^i.1`$ is degree zero (it acts by differentiation) which is then killed by $`\mathrm{d}`$. Therefore the only term that could contribute in the first line, for example, is from $`Q_2^2`$ i.e. $`(e^ie^jv)\mathrm{d}e_ie_jw`$ which is zero because of the second differentiation on $`v`$. The difference between the two expressions is of course $$v\mathrm{d}w\mathrm{d}wv=\frac{\mathrm{}}{2}\mathrm{d}[v,w]+O(\mathrm{}^3),$$ i.e. the Poisson-compatible preconnection as it should be. The example of ‘noncommutative spacetime’ with nonassociative differentials $$t\mathrm{d}x_i\mathrm{d}x_it=\mathrm{d}tx_ix_i\mathrm{d}t=\frac{\mathrm{}}{2}\mathrm{d}x_i$$ is very different from the associative (but not canonical) differential calculus usually used for this model. It represents a different approach that may overcome some of the structural problems encountered previously (such as to find the canonical Dirac operator, see below). Note that in the physical application the deformation parameter that we have denoted $`\mathrm{}`$ should be denoted by another symbol and is expected to be of the order of the Planck time $`10^{44}`$s. Finally, the Mackey quantisaition case is more complicated but from the curvature computations in Section 5 we conclude again that the natural deformed calculus $`\mathrm{\Omega }(C^{\mathrm{}}(N)>U_{\mathrm{}}(𝔤))`$ is nonassociative at least for simple Lie algebras, because the CBH part is. It is interesting to note that the curvature comes from this part alone. ### 8.3. Isospectral quantum Dirac operator Here we conclude with an example to demonstrate that the categorical deformation method outlined in Section 2 is very powerful indeed and quantizes almost any natural construction. In other words, when a quantisation is expressed as a cochain module algebra twist this has great consequences. Specifically, in another approach to noncommutative geometry it is normal to look for an analogue of the Dirac operator in the form of a ’spectral triple’ \[C\] obeying some axioms. These axioms are natural from an associative point of view but it is well known that important examples such as $`_q[G]`$ do not admit operators obeying exactly those axioms. We see by contrast that there is a natural deformation of any classical Dirac operator on the classical phase space but it will obey a variation of Connes axioms due to the hidden nonassociavity in the underlying differential calculus and elsewhere. We explain now that such an approach agrees with recent ’isospectral deformation’ proposal for the Dirac operator on $`_q[SU_2]`$ in \[DLSSV\]. On the other hand, it is more categorical and works in principle for all quantum groups $`_q[G]`$, and moreover works for our more conventional quantisations such as $`U_{\mathrm{}}(g)`$ and the Mackey quantisation to provide (in principle at least) some type of Dirac operators on them. We consider for the sake of discussion only the case where the classical and hence quantum cotangent and spin bundles are trivial so that the spin bundle in particular has the form $`VA`$ where $`A=C^{\mathrm{}}(M)`$ and $`V`$ is ostensibly the representation space for the spin group. We do require everything to be covariant under a background Lie algebra $``$ (or Hopf algebra $`H`$) to induce the quantisation given a cochain. This is not a problem in the case $`M=G`$ a Lie group (a covariant Dirac operator). The short version of the quantisation is then as follows: we consider the classical Dirac operator $`D:VAVA`$ and to this we apply the functor $`𝒯`$ in Section 2 to obtain a map $`𝒯(D):𝒯(VA)𝒯(VA)`$. As in Section 2 we have to allow that although $`𝒯`$ acts as the identity on objects and morphisms (so $`𝒯(A)=A`$ which becomes the deformed algebra $`A_F`$, $`𝒯(V)=V`$, $`𝒯(D)=D`$), it is nontrivial as a monoidal functor and in the sense of potentially nontrivial natural isomorphisms $`𝒯(V)𝒯(A)𝒯(VA)`$ with certain properties in relation to $``$. We refer to \[M2\] for an introduction. Here these isomorphisms are given by the action of $`F^1`$. Putting these facts together, we have the deformed Dirac operator: $$D_{}:A_FVA_FV,D_{}(av)=FD(F^1(av)).$$ The main thing to note about this construction is that since it is given by conjugation by $`F`$ as an operator, it does not change the spectrum in the Hilbert space setting. It should be remarked that one would still need a lot of analysis to make these remarks fully precise. The only other subtlety is to identify $`A_F`$ (which is the same vector space as $`A`$ with the deformed product) explicitly as $`_q[G]`$ in the case $`A=[G]`$. This is not trivial but we note that $`𝒯`$ respects sums so if one has made a Peter-Weyl decomposition of $`A`$ into a direct sum of matrix algebras (as is possible for $`A=[G]`$ for simple $`G`$) and likewise decompose $`_q[G]`$ in its Peter-Weyl decomposition, we can identify the summands as matrix coalgebras. This is our interpretation of the proposed Dirac operator in \[DLSSV\]. On the other hand, $`D_{}`$ lives in a nontrivial monoidal category and has properties in which the non-associativity of the category will surely play a role. It is known that the axioms in \[C\] are not satisfied and we would propose to replace them by ones that take this hidden nonassociavity into account. Finally, the longer answer to the deformed Dirac operator here is to ’get inside’ its construction. One can do this too in principle as we outline now. Thus, we break $`D`$ into a series of morphisms all covariant under our background Hopf algebra. We also suppose for the sake of discussion that $`M`$ is parallelizable (e.g. $`M`$ a Lie group) so that its (covariant) differential calculus has the form $`\mathrm{\Omega }^1(M)=A\mathrm{\Lambda }^1`$ as a (trivial) bundle associated to $`\mathrm{\Lambda }^1`$. We write $`\mathrm{d}a=(^ia)\tau _i`$ where $`\tau _i`$ are a basis of $`\mathrm{\Lambda }^1`$ and take this as a definition of the partial derivatives. Finally, we assume ‘$`\gamma `$-matrices’ of some form $`\gamma :\mathrm{\Lambda }^1VV`$ so that $$D(av)=^i(a)\gamma _i(v)=(\mathrm{id}\gamma )(\mathrm{d}\mathrm{id})(av)$$ expresses $`D`$ as a composition of morphisms (here $`\gamma _i=\gamma (\tau _i())`$ would be the more conventional point of view). Note that the defining relations among the $`\gamma `$ also needs to be invariant under the background symmetry, e.g. by an invariant metric. We now define $`\gamma _{}(\tau v)=\gamma (F^1(\tau v))`$ for all $`\tau \mathrm{\Lambda }^1`$, $`vV`$ by the same reasoning as above, i.e. the functor $`𝒯`$. Likewise we have $`\mathrm{d}_{}=F\mathrm{d}`$ when $`\mathrm{d}:AA\mathrm{\Lambda }^1`$. Note that in the above we have not deformed $`\mathrm{d}`$ and indeed this is not deformed if we consider it to $`\mathrm{\Omega }^1`$ and identify $`𝒯(\mathrm{\Omega }^1)=\mathrm{\Omega }^1`$ in the deformed theory, $`\mathrm{d}_{}`$ is a slightly different object. This now expresses the above as $$D_{}=FDF^1=(\mathrm{id}\gamma _{})\mathrm{\Phi }_{A,\mathrm{\Lambda }^1,V}(\mathrm{d}_{}\mathrm{id})$$ in view of the diagram: $$\begin{array}{ccccc}𝒯(AV)& & 𝒯(A)𝒯(V)& & \\ \mathrm{d}& & \mathrm{d}& \mathrm{d}_{}& \\ 𝒯((A\mathrm{\Lambda }^1)V)& & 𝒯(A\mathrm{\Lambda }^1)𝒯(V)& & ((𝒯(A)𝒯(\mathrm{\Lambda }^1))𝒯(V)\\ \|\|& & & & \mathrm{\Phi }_{A,\mathrm{\Lambda }^1,V}\\ 𝒯(A(\mathrm{\Lambda }^1V))& & 𝒯(A)𝒯(\mathrm{\Lambda }^1V)& & (𝒯(A)(𝒯(\mathrm{\Lambda }^1)𝒯(V))\\ \gamma & & \gamma & \gamma _{}& \\ 𝒯(AV)& & 𝒯(A)𝒯(V)& \end{array}$$ The large middle cell here commutes by definition of a monoidal functor (the associator in the initial category of $`H`$-covariant objects is trivial). The upper and lower left cells commute because of operations on different spaces. The upper and lower right cells commute by the definitions of $`\mathrm{d}_{}`$ and $`\gamma _{}`$ respectively. The horizontal arrows are all given by the action of $`F`$. The vertical composition on the left is $`𝒯(D)=D`$, while the vertical composition on the right is the definition of $`D_{}`$. As to the $`\gamma _{}`$, one should write the defining relations of $`\gamma `$ as commuting diagrams, apply the functor $`𝒯`$ to obtain commuting diagrams in the deformed category, and use $`F`$ to interpret them as $`\gamma _{}`$ relations in a similar manner to the above. The result is, for example, a deformed set of Clifford relations involving now $`\mathrm{\Phi }`$ and the deformation of the flip map induced by $`F`$ (as a symmetry in the category). The actual relations would depend on the classical set up which need not be the usual classical Clifford relations if the frame group is not the usual one (e.g. one may use a Lie group to frame itself).
warning/0506/hep-th0506064.html	ar5iv	text	# Contents ## Chapter 1 An introduction to non(anti)commutative geometry and superstring theory ### 1.1 A brief introduction to noncommutative and non(anti)commutative field theory #### 1.1.1 The Moyal product ##### Weyl transform definition Noncommutative geometry deals with manifolds whose coordinates do not commute. This kind of manifold appeared in physics much before noncommutative geometry itself was born as a branch of mathematics . A well-known example is quantum phase space. This is a manifold described by $`2n`$ operator coordinates $`\widehat{X}_1`$,…,$`\widehat{X}_n`$,$`\widehat{P}_1`$,…,$`\widehat{P}_n`$, satisfying nontrivial commutation relations $`[\widehat{X}_i,\widehat{P}_j]=i\mathrm{}\delta _{ij}`$ (1.1) $`[\widehat{X}_i,\widehat{X}_j]=[\widehat{P}_i,\widehat{P}_j]=0`$ (1.2) It is expected that the geometric nature of spacetime will be modified at very short distances. The physical idea underlying modern noncommutative geometry is that a quantum spacetime will be uncovered then, where the usual trivial commutation relations among coordinates are no longer valid and noncommutativity emerges as $$[\widehat{X}^\mu ,\widehat{X}^\nu ]=i\theta ^{\mu \nu }\left(\widehat{X}\right)$$ (1.3) with $`\theta ^{\mu \nu }=\theta ^{\nu \mu }`$. In the limit $`\theta ^{\mu \nu }0`$ ordinary commutative geometry must be recovered. Coordinate algebra (1.3) gives rise to the spacetime uncertainty relations $$\mathrm{\Delta }X^\mu \mathrm{\Delta }X^\nu \frac{1}{2}\|\theta ^{\mu \nu }\|$$ (1.4) In (1.3) the possibility is left open that time may be involved in noncommutativity. In this case (1.3) would not be simply a generalization of the quantum mechanical commutation relations (1.2) and it is somehow expected that it may clash with quantum mechanics. For a while I will not worry about this, I will reconsider the case of noncommuting time later on. The connection between noncommutative geometry and quantum mechanics is easily seen when the latter is discussed in the Weyl formalism. In this formalism an explicit map between functions $`f(x,p)`$ of the phase space variables $`x`$, $`p`$ and corresponding operators $`\widehat{O}_f(\widehat{X},\widehat{P})`$ is constructed, where $`\widehat{X}`$, $`\widehat{P}`$ are noncommuting operators corresponding to classical variables $`x`$, $`p`$. I will briefly discuss this formalism, in the case of two variables $`x^1`$, $`x^2`$, with corresponding operators satisfying $`[\widehat{X}^1,\widehat{X}^2]=\theta ^{12}`$ with constant $`\theta `$. After this discussion the natural embedding of noncommutative geometry in quantum mechanics will be clear to the reader . On “phase space” described by coordinates $`x_1,x_2`$ let us consider a function $`f(x^1,x^2)`$. Given its Fourier transform $$\stackrel{~}{f}(\alpha _1,\alpha _2)=d^2xe^{i(\alpha _1x^1+\alpha _2x^2)}f(x^1,x^2)$$ (1.5) we can define an operator $`\widehat{O}_f(\widehat{X^1},\widehat{X^2})`$ as follows $$\widehat{O}_f(\widehat{X^1},\widehat{X^2})=\frac{1}{(2\pi )^2}d^2\alpha U(\alpha _1,\alpha _2)\stackrel{~}{f}(\alpha _1,\alpha _2)$$ (1.6) where $$U(\alpha _1,\alpha _2)=e^{i(\alpha _1\widehat{X^1}+\alpha _2\widehat{X^2})}$$ (1.7) Making use of Baker-Campbell-Hausdorff formula we find $$U(\alpha _1,\alpha _2)U(\beta _1,\beta _2)=e^{\frac{i}{2}(\alpha _1\beta _2\alpha _2\beta _1)\theta ^{12}}U(\alpha _1+\beta _1,\alpha _2+\beta _2)$$ (1.8) The map $`f\widehat{O}_f`$ defines the Weyl-Moyal correspondence between functions on phase space and operators. Now we would like to determine which function corresponds to the operator $`\widehat{O}_f\widehat{O}_g`$. We know that in general $`\widehat{O}_f\widehat{O}_g\widehat{O}_g\widehat{O}_f`$, so we expect some noncommutative deformation of the ordinary product to arise by the Weyl-Moyal correspondence. $`\widehat{O}_f\widehat{O}_g={\displaystyle \frac{1}{(2\pi )^4}}{\displaystyle d^2\alpha d^2\beta U(\alpha _1,\alpha _2)U(\beta _1,\beta _2)\stackrel{~}{f}(\alpha _1,\alpha _2)\stackrel{~}{g}(\beta _1,\beta _2)}=`$ (1.9) $`={\displaystyle \frac{1}{(2\pi )^4}}{\displaystyle d^2\alpha d^2\beta U(\alpha _1+\beta _1,\alpha _2+\beta _2)e^{\frac{i}{2}(\alpha _1\beta _2\alpha _2\beta _1)\theta ^{12}}\stackrel{~}{f}(\alpha _1,\alpha _2)\stackrel{~}{g}(\beta _1,\beta _2)}`$ (1.10) (1.11) By making the change of variables $`\gamma _1=\alpha _1+\beta _1`$, $`\delta _1=\frac{1}{2}(\alpha _1\beta _1)`$, $`\gamma _2=\alpha _2+\beta _2`$, $`\delta _2=\frac{1}{2}(\alpha _2\beta _2)`$ we obtain $`\widehat{O}_f\widehat{O}_g={\displaystyle \frac{1}{(2\pi )^4}}{\displaystyle }d^2\gamma d^2\delta U(\gamma _1,\gamma _2)e^{\frac{i}{2}\theta ^{12}(\gamma _1\delta _2\delta _1\gamma _2)}\stackrel{~}{f}({\displaystyle \frac{\gamma _1}{2}}+\delta _1,{\displaystyle \frac{\gamma _2}{2}}+\delta _2)`$ (1.12) $`\stackrel{~}{g}({\displaystyle \frac{\gamma _1}{2}}\delta _1,{\displaystyle \frac{\gamma _2}{2}}\delta _2)`$ (1.13) Let us define Moyal product between two functions $`f`$ and $`g`$ in $`R^{2n}`$ as $$\left(fg\right)(x)=e^{\frac{i}{2}\theta ^{ij}_i_j^{}}f(x)g(x^{})\|_{x=x^{}}$$ (1.14) In our simple two-dimensional case it becomes $$(fg)(x)=e^{\frac{i}{2}\theta ^{12}(_1_2^{}_2_1^{})}f(x)g(x^{})\|_{x=x^{}}$$ (1.15) In momentum space we can obtain the following formula for the Fourier transform of $`fg`$ $$\stackrel{~}{fg}(\gamma _1,\gamma _2)=\frac{1}{(2\pi )^2}d^2\delta e^{\frac{i}{2}\theta ^{12}(\gamma _1\delta _2\gamma _2\delta _1)}\stackrel{~}{f}(\frac{\gamma _1}{2}+\delta _1,\frac{\gamma _2}{2}+\delta _2)\stackrel{~}{g}(\frac{\gamma _1}{2}\delta _1,\frac{\gamma _2}{2}\delta _2)$$ (1.16) From this and (1.13) it is clear that $$\widehat{O}_f\widehat{O}_g=\frac{1}{(2\pi )^2}d^2\gamma U(\gamma _1,\gamma _2)\stackrel{~}{fg}(\gamma _1,\gamma _2)=\widehat{O}_{fg}$$ (1.17) So Moyal product (1.14) naturally emerges in the context of quantum mechanics, when the latter is expressed in the Weyl-Moyal formalism. It is the functional product corresponding to operator product between quantum observables. Applying (1.14) to the special case $`f=x^i`$, $`g=x^j`$, we obtain the coordinate algebra $$x^ix^jx^jx^i=[x^i,x^j]_{}=i\theta ^{ij}$$ (1.18) Therefore the quantum commutation relations we began with are reproduced in the functional formalism as $``$ commutators. ##### Translation covariance and associativity as a definition Moyal product (1.14) can also be obtained from a general discussion concerning the algebraic requirement of associativity and the geometric requirement of covariance with respect to translations. These two properties uniquely determine Moyal product. Before I discuss this, I will introduce the general ideas concerning a field called deformation quantization and its connections with modern noncommutative geometry. Ordinary geometry is based on the concept of point. This is not true anymore for noncommutative geometry, since a noncommutative manifold is completely defined in terms of the properties of the algebra of functions on it . In ordinary geometry many sets of points can be completely described when the algebra $`A`$ of functions on them with values in $`R`$ or $`C`$ is known. A finite dimensional vector space $`V`$ is a familiar example of this, since the space of functions $`f:VR`$ (or $`C`$) is the dual space $`V^{}`$, which is isomorphic to $`V`$. In this case studying the algebra of functions on the manifold or the manifold itself is the same thing. We can consider the more general case of a $`C^{}`$ algebra $`A`$, i.e. an algebra endowed with a norm and an involution. Every $`C^{}`$ algebra is isomorphic to the algebra $`A^{}`$ of complex continuos functions on a certain compact space $`V`$. When $`A^{}`$ is commutative we can go back to the space $`V`$, that can be described as a set of points in ordinary geometry. When $`A^{}`$ is noncommutative, instead, going back to the space $`V`$ can be very complicated and in some cases impossible. However, this is irrelevant for the purpose of studying a physical theory, since all the needed information are encoded in $`A^{}`$. A recipe to obtain a theory on noncommutative space from a given one on ordinary space is the following. Consider the algebra of functions with values in R (or C), deform its product to a new, noncommutative one, that I will call $``$, defined in terms of a parameter $`\mathrm{}`$. In the limit $`\mathrm{}0`$ one must recover ordinary, commutative case. Now rewrite the old theory replacing all ordinary products with $``$ products, and think of the new theory as a deformation of the ordinary one, defined on noncommutative space. This I will call the natural deformation of a theory. It is not the only possible definition of a noncommutative generalization of an ordinary theory and I will discuss this point in more detail in section 1.1.3. Given two functions $`f`$, $`g`$, I will denote their ordinary, commutative product as $`fg`$. I will deform it in the following way $$fgfg+\mathrm{}P(f,g)+𝒪(\mathrm{}^2)$$ (1.19) where $`P(f,g)`$ is a bilinear operator in the two functions $`f`$, $`g`$ and $`\mathrm{}`$ is the parameter governing noncommutativity. The example of quantum phase space discussed before suggests a good candidate for the bilinear operator $`P`$. When the manifold we are considering is endowed with a Poisson structure $`\{,\}_P`$, we will choose $$P(f,g)=\{f,g\}_P=P^{\mu \nu }_\mu _\nu ^{}f(x)g(x^{})\|_{x=x^{}}=f\stackrel{}{}_\mu P^{\mu \nu }\stackrel{}{}_\nu g$$ (1.20) (The last equality is just to present a different notation. It will be preferred since it is more suitable to superspace extension, where the presence of fermionic indices makes different orderings inequivalent). In the 70’s the deformation of Poisson manifolds was studied in a completely different context. In the paper by Bayen et al. a different approach to quantization was proposed. Quantization had to be understood as a deformation of the algebraic structure of functions and not as a radical change in the nature of physical observables. Moyal product $``$ was then introduced with the goal of reinterpreting quantum mechanics in the context of algebraic deformations. In particular, in an analysis of the possible noncommutative but associative products that can be obtained as a perturbative series in the parameter $`\mathrm{}`$ was performed. The results obtained there are very interesting when they are reread in the light of the new ideas of noncommutative geometry. Consider a manifold $`\mathrm{\Omega }`$ endowed with a Poisson structure $`P`$, where a set of derivatives $`_\mu `$ is defined such that $`_\mu P=0`$. We will also assume that this set of derivatives is torsion free and without curvature. We define a generic product $``$ on $`\mathrm{\Omega }`$ by the smooth function $$u(z)=\underset{r=0}{\overset{\mathrm{}}{}}a_r\left(\frac{z^r}{r!}\right)$$ (1.21) with $`a_0=a_1=1`$, as $$fg=\underset{r=0}{\overset{\mathrm{}}{}}\mathrm{}^r\frac{a_r}{r!}P^r(f,g)$$ (1.22) where $$P^r(f,g)=P^{\mu _1\nu _1}\mathrm{}\mathrm{}.P^{\mu _r\nu _r}_{\mu _1}\mathrm{}_{\mu _r}f_{\nu _1}\mathrm{}_{\nu _r}g$$ (1.23) One can show that the exponential function is the only possible choice for $`u`$ leading to an associative product, i. e. satisfying $$(fg)h=f(gh)$$ (1.24) To show this one imposes (1.24), writing every $``$ product explicitly as in (1.22, 1.23). Order by order in $`\mathrm{}`$ one gets constraints on the coefficients $`a_r`$. The proof makes a strong use of the assumptions on the derivative $``$, since one needs to exchange derivatives and to pass the Poisson tensor $`P^{\mu \nu }`$ through derivatives without getting extra terms from commutators. Finally one obtains that (1.24) is satisfied if and only if $`a_r=1r`$ and this uniquely identifies the function $`u`$ with the exponential. Summarizing, under the assumptions made for $`\mathrm{\Omega }`$, $``$ and $`P`$, the unique associative $``$ product has the form: $$fg=e^{(\mathrm{}P)}(f,g)$$ (1.25) (modulo a constant overall factor and linear changes of variables). If at least one of the three hypotesis is not satisfied ($`P`$ constant with respect to $``$, $``$ without torsion and curvature), then Moyal product is not associative anymore. This can be easily seen by considering second and third order terms in $`\mathrm{}`$. Once the product $``$ is known, the commutation relations among coordinates are determined by considering the special case of two coordinates themselves as functions $`f`$ and $`g`$ . If $`\mathrm{\Omega }`$ is flat spacetime described by coordinates $`\left\{x^\mu \right\}`$ and ordinary derivatives, we can take as a Poisson structure the one associated to a constant antisymmetric matrix $`P^{\mu \nu }`$. In this case we obtain the coordinate algebra $$[x^\mu ,x^\nu ]_{}x^\mu x^\nu x^\nu x^\mu =2\mathrm{}P^{\mu \nu }.$$ (1.26) Usually in the definition of the commutator an $`i`$ is factorized so that the matrix $`P`$ is hermitian. Moreover, the parameter $`\mathrm{}`$ is sometimes absorbed into the definition of the matrix $`P`$. In flat spacetime the choice of a constant $`P^{\mu \nu }`$ is deeply related to translation invariance. In fact, if we want to deform a theory with this symmetry, the only deformation of the coordinate algebra that preserves it is the one associated to a constant symplectic matrix. Consider the commutation relations $`[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }(x)`$ that we would like to implement in a certain theory originally defined in terms of commuting coordinates $`x^\mu `$. Suppose the original theory to be symmetric with respect to the transformation $`xx^{}`$. For the symmetry to be preserved in the deformed theory the new coordinate algebra must be invariant with respect to that trasformation. When we say invariant we mean that the functional dependence on $`x`$ variable must not change under the transformation, that is $`[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }(x)_{xx^{}}[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }(x^{})`$ (1.27) Note that the matrix $`\theta ^{\mu \nu }`$ only transforms punctually and does not become a new matrix $`\theta ^{\mu \nu }`$. The matrix $`\theta ^{\mu \nu }`$ is arbitrarily chosen, defines the noncommutative manifold and must be same for all $`x`$ on the manifold. Again, let us consider the case of flat spacetime. We would like to deform a Poincaré invariant theory. We will first consider translations $`xx+a`$ to see which conditions must be imposed on the matrix $`\theta ^{\mu \nu }`$ for the deformed algebra not to break this symmetry. $`[x^\mu ,x^\nu ]=[x^\mu +a^\mu ,x^\nu +a^\nu ]=[x^\mu ,x^\nu ]`$ (1.28) So $`\theta `$ must satisfy the constraint $$\theta ^{\mu \nu }(x+a)=\theta ^{\mu \nu }(x)$$ (1.29) Since $`\theta ^{\mu \nu }`$ must be a local function, it has to be constant. Therefore, in flat spacetime the only nontrivial deformation preserving translation invariance is the one with constant commutators. Now we will consider Lorentz invariance . Two different kinds of Lorentz transformations can be considered, the ones where the observer moves while the particle stands still (“observer” Lorentz transformations) and the ones where the particle is boosted or rotated and the observer is fixed (“particle” Lorentz transformations). In the first case it is sufficient for the physics of the system not to change that the matrix $`\theta ^{\mu \nu }`$ transforms covariantly. In the second case instead the matrix $`\theta ^{\mu \nu }`$ must not transform, since the coordinate algebra must remain unaltered while moving from $`x`$ to $`x^{}`$. Thus in this case physics changes under the transformation and the symmetry is broken. Let us explicitly consider the “particle” Lorentz transformation $`x^\mu x^\mu =\mathrm{\Lambda }_\nu ^\mu x^\nu `$. It happens that $$[x^\mu ,x^\nu ]=[\mathrm{\Lambda }_\rho ^\mu x^\rho ,\mathrm{\Lambda }_\sigma ^\nu x^\sigma ]=\mathrm{\Lambda }_\rho ^\mu [x^\rho ,x^\sigma ]\mathrm{\Lambda }_\sigma ^\nu =\mathrm{\Lambda }_\rho ^\mu \theta ^{\rho \sigma }\mathrm{\Lambda }_\sigma ^\nu _{D>2}\theta ^{\mu \nu }$$ (1.30) We conclude that noncommutative theories in $`D>2`$ dimensions cannot preserve “particle” Lorentz transformation, while “observer” Lorentz transformation are not broken by the deformation. An exception to this general rule is the two-dimensional case, where every antisymmetric matrix is a number times the Ricci tensor $`ϵ^{\mu \nu }`$, which is Lorentz invariant. In most papers concerning noncommutative field theory the following choice is made $`\theta ^{0i}=0,\theta ^{ij}0`$ (1.31) Time is “isolated” with respect to spacial directions and Lorentz symmetry is manifestly broken. As anticipated in the beginning of this section, time-space noncommutativity is likely to cause a breakdown of the usual framework of quantum mechanics. Actually, it has been shown that time-space noncommutativity is responsible for unitarity and causality problems in noncommutative field theory (see section 1.1.2). To avoid this, the restriction (1.31) is applied in most work concerning noncommutative field theory. Finally I would like to point out that the discussion about symmetries I have presented here is based on the assumption that the symmetry group is undeformed (i.e. it is a classical symmetry group and not a quantum group). This means that parameters of symmetry transformations are commuting. This is not the only possible way to proceed. There is a branch of mathematics called Quantum Algebra that studies the deformation of symmetry groups. An interesting example is the $`\kappa `$-deformation of Minkowski space , where parameter and coordinate algebras have an identical structure. Since Minkowski spacetime can be defined as the quotient between Poincaré and Lorentz groups, it may seem natural to take also into consideration deformations of Minkowski space that are accompanied by an analogous deformation in the translation symmetry group. I briefly summarize the results obtained, in the special case $`\mathrm{\Omega }=R^{2n}`$ (the extension to Minkowski signature is straightforward). The only product $``$ defined as in (1.19, 1.20) that is associative and that preserves translation invariance is Moyal product $$(fg)(x)=e^{\frac{i}{2}\mathrm{}\theta ^{ij}_i_j^{}}f(x)g(x^{})\|_{x=x^{}}$$ (1.32) that generates the coordinate algebra $$[x^i,x^j]_{}=i\mathrm{}\theta ^{ij}$$ (1.33) As we have seen before, Moyal product is also naturally obtained in quantum mechanics through the Weyl-Moyal correspondence defined between quantum operators and functions on “phase-space”. ##### Properties of Moyal product Here I will summarize some useful properties of Moyal product $``$. Associativity and covariance with respect to translations have been already discussed. 1. The $``$ product between exponential functions reflects from the functional point of view the well-known Baker-Campbell-Hausdorff formula $`e^{ikx}e^{iqx}=e^{i(k+q)x}e^{\frac{i}{2}(k\theta q)}`$ (1.34) $`k\theta qk^\mu q^\nu \theta _{\mu \nu }`$ (1.35) 2. By making use of the previous formula we can obtain the representation of $``$ in momentum space $$\left(fg\right)(x)=\frac{1}{(2\pi )^8}d^4kd^4q\stackrel{~}{f}(k)\stackrel{~}{g}(q)e^{\frac{i}{2}(k\theta q)}e^{i(k+q)x}$$ (1.36) 3. Commutativity is recovered under integration $$(fg)(x)d^4x=(gf)(x)d^4x=(fg)(x)d^4x$$ (1.37) since all the corrections in (1.32) with respect to the ordinary product are total derivatives, because of the antisymmetry of $`\theta ^{\mu \nu }`$. 4. A cyclicity property can be deduced from the previous relation $$(f_1f_2\mathrm{}f_n)(x)d^4x=(f_nf_1\mathrm{}f_{n1})(x)d^4x$$ (1.38) 5. Finally, $``$ has the following behavior with respect to complex conjugation $$\left(fg\right)^{}=g^{}f^{}$$ (1.39) because of the antisymmetry of $`\theta ^{\mu \nu }`$. Clearly $`ff`$ is real when $`f`$ is real, but when both $`f`$ and $`g`$ are real, $`fg`$ is generally complex. #### 1.1.2 The natural Moyal deformation of a field theory In this section I will discuss the main properties of noncommutative field theories that are obtained from ordinary ones by replacing ordinary products with Moyal $``$ products (1.32) in the action. This is what I will call the natural deformation of a given field theory. I will first discuss the simple case of scalar field theory with $`\mathrm{\Phi }^4`$ interaction . This is chosen for simplicity and most of the features we will find in this case can be easily generalized to more complicated situations. I will then move to gauge theories, to see how gauge invariance is modified in noncommutative space. In this first two subsections I will only take into consideration the restricted case (1.31). In the last subsection I will instead discuss unitarity and causality problems arising when time is involved in noncommutativity. ##### A simple example: The scalar $`\mathrm{\Phi }^4`$ theory Let us consider the natural noncommutative deformation of a given ordinary field theory, for instance the scalar theory with $`\mathrm{\Phi }^4`$ interaction. We have already seen that we can obtain the deformed theory by replacing ordinary products with $``$ products everywhere in the action. We choose Moyal product because we want to preserve translation invariance in the deformed theory and we want associativity. The action for the noncommutative theory is $$S\left[\mathrm{\Phi }\right]=d^4x\left[\frac{1}{2}_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }\frac{m^2}{2}\mathrm{\Phi }\mathrm{\Phi }\frac{\lambda }{4!}\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\right]$$ (1.40) Property (1.37) implies that the quadratic part of the action does not receive corrections from the star products. Only the interaction term is modified, so the free theory is the same as the ordinary one. The noncommutative theory is built on the same Fock space as the commutative one, but it has different interactions. This feature is common to all theories obtained as deformation of ordinary ones by implementing Moyal product, since it just relies on property (1.37). We can easily deduce Feynman rules from (1.36). Introducing the Fourier components $`\varphi (k)`$ of $`\mathrm{\Phi }(x)`$ $$\mathrm{\Phi }(x)=\frac{1}{(2\pi )^4}d^4ke^{ikx}\varphi (k)$$ (1.41) We obtain $`S_{\mathrm{int}}={\displaystyle \frac{\lambda }{4!}}{\displaystyle d^4x\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }}`$ (1.42) $`={\displaystyle \frac{1}{(2\pi )^{16}}}{\displaystyle \frac{\lambda }{34!}}{\displaystyle d^4k_1\mathrm{}d^4k_4\varphi (k_1)\varphi (k_2)\varphi (k_3)\varphi (k_4)(2\pi )^4\delta ^{(4)}(\underset{i=1}{\overset{4}{}}k_i)}`$ (1.43) $`\left[\mathrm{cos}{\displaystyle \frac{k_1\theta k_2}{2}}\mathrm{cos}{\displaystyle \frac{k_3\theta k_4}{2}}+\mathrm{cos}{\displaystyle \frac{k_1\theta k_3}{2}}\mathrm{cos}{\displaystyle \frac{k_2\theta k_4}{2}}+\mathrm{cos}{\displaystyle \frac{k_1\theta k_4}{2}}\mathrm{cos}{\displaystyle \frac{k_2\theta k_3}{2}}\right]`$ (1.44) (1.45) We conclude that the only difference between the natural deformation of a field theory and the field theory itself is a phase factor depending on momenta and noncommutativity parameter $`\theta `$, appearing in front of every vertex in the Feynman rules. This procedure can be clearly generalized to other field theories. Now I’m going to discuss how this phases modify perturbation theory, in particular ultraviolet behaviour and renormalization. Since the phases appearing in front of vertices depend on the momenta, when we compute the contribution coming from a certain diagram we have to distinguish between two different situations. If the phase is only depending on external momenta, it does not affect loop integrations and thus it does not modify the degree of divergence. This case we will call planar. Instead, when the phase factor depends on internal, loop momenta, it generally modifies the ultraviolet behavior of the diagram. This case we will call nonplanar. So a single diagram in the ordinary theory decomposes in various planar and nonplanar contributions, depending on the ordering of momenta in the vertices. A nice feature of natural Moyal deformations of ordinary field theories is that nonplanar graphs always display a better ultraviolet behavior with respect to the corresponding planar ones, since the phase acts as a regulator. So one can say that such deformation of a renormalizable theory will also be renormalizable. It will display the same degree of divergence in planar diagrams and a lower degree of divergence in nonplanar ones . A general discussion of renormalizability properties of noncommutative field theory can be found in . Now I would like to discuss a typical feature of noncommutative field theories called UV/IR mixing. To this purpose I will present the result of the 1-loop computation for the renormalized two-point function $`\mathrm{\Gamma }_{\mathrm{ren}}^{(2)}`$ in the case of $`\mathrm{\Phi }^4`$ theory. I will not give any detail about the computation. The interested reader should refer to . Let $`\lambda `$ be the coupling constant of the theory, $`\mathrm{\Lambda }`$ the ultraviolet cutoff, $`M`$ the renormalized mass, $`p`$ the incoming momentum. Moreover we will define $`pkp\theta \theta k=p_\mu \theta ^{\mu \rho }\theta _\rho ^\nu k_\nu `$, where $`\theta ^{\mu \nu }`$ is the noncommutativity matrix characterizing the theory. One finds that the renormalized $`\mathrm{\Gamma }^{(2)}`$ in the limit $`\mathrm{\Lambda }_{\mathrm{eff}}\frac{1}{\left(pp+\frac{1}{\mathrm{\Lambda }^2}\right)^{\frac{1}{2}}}0`$ takes the form $$\mathrm{\Gamma }_{\mathrm{ren}}^{(2)}(p,M,\mathrm{\Lambda })=p^2+M^2+\frac{\lambda }{96(2\pi )^2(pp+\frac{1}{\mathrm{\Lambda }^2})}\frac{\lambda M^2}{96\pi ^2}\mathrm{ln}\frac{1}{M^2(pp+\frac{1}{\mathrm{\Lambda }^2})}+𝒪(\lambda ^2)$$ (1.46) If we then take the limit $`\mathrm{\Lambda }\mathrm{}`$ we observe that an infrared divergence appears when we take $`p0`$. If we instead take the limit $`p0`$ first, we discover that the cutoff does not appear explicitly anymore and the two-point function diverges for $`\mathrm{\Lambda }\mathrm{}`$. So we observe an interesting connection between the UV and IR behaviors in the extra terms appearing in the two-point function because of noncommutativity. This is known in the literature as UV/IR mixing. Finally, we will consider the limit $`\theta 0`$. In this limit we expect to obtain the standard result for the renormalized $`\mathrm{\Gamma }^2`$ of ordinary $`\mathrm{\Phi }^4`$ theory. We have already discussed the fact that a diagram in the ordinary theory splits in planar and nonplanar parts in the noncommutative theory. In (1.46) the subtraction made to obtain the renormalized mass only took into consideration the divergences coming from the planar graphs. The last two terms represent the contribution coming from nonplanar diagrams. In the limit $`\theta 0`$ one may verify that by adding to the planar contributions the nonplanar ones in the computation of the renormalized mass, one obtains the well-known result for ordinary $`\mathrm{\Phi }^4`$. ##### Gauge theories Up to now I have considered scalar theory for simplicity. Now I would like to discuss Yang-Mills theories. In ordinary geometry these theories are constructed by promoting a global invariance to a local one. In general the gauge group is nonabelian. First of all a remark is needed regarding the choice of the gauge group. In noncommutative geometry described by Moyal product it is easy to show that $`SU(n)`$, $`SO(n)`$ and $`Sp(n)`$ are not closed any more and the same is valid for the corresponding Lie algebras . This is due to the fact that a nontrivial trace part appears in the product of two traceless matrices. So we will consider $`U(n)`$ gauge theory as our example<sup>1</sup><sup>1</sup>1The cases $`SO(n)`$, $`Sp(n)`$ seem to allowed from the subtle string theory discussion in . The ordinary gauge theory is described by an $`n\times n`$ hermitian matrix of vectors, $`A_\mu `$. This transforms as follows under the local gauge symmetry with parameter $`\lambda `$ (also an $`n\times n`$ matrix) $$\delta _\lambda A_\mu =_\mu \lambda +i[\lambda ,A_\mu ]$$ (1.47) The field strength $`F_{\mu \nu }`$ and its transformation law are given by $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu i[A_\mu ,A_\nu ]`$ (1.48) $`\delta _\lambda F_{\mu \nu }=i[\lambda ,F_{\mu \nu }]`$ (1.49) The pure gauge action is given by $$Tr\left(F^{\mu \nu }F_{\mu \nu }\right)$$ (1.50) where the trace acts on gauge indices and the integral is taken over spacetime variables. The invariance of (1.50) under (1.49) follows from cyclicity of the trace. We are going to construct the natural Moyal deformation of the theory by substituting ordinary products with $``$ products in the action. This time we have matrix-valued fields, so the product between two of them will be the tensor product between $``$ and the matrix product. We obtain the action $$Tr\left(F^{\mu \nu }F_{\mu \nu }\right)$$ (1.51) If we also substitute ordinary products with Moyal products in the gauge transformation and definition of the field strength $`\delta _\lambda A_\mu =_\mu \lambda +i\lambda A_\mu iA_\mu \lambda `$ (1.52) $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu iA_\mu A_\nu +iA_\nu A_\mu `$ (1.53) $`\delta _\lambda F_{\mu \nu }=i\lambda F_{\mu \nu }iF_{\mu \nu }\lambda `$ (1.54) we find that the gauge transformation is still a symmetry of the action. In the noncommutative case the invariance of (1.51) under (1.54) is more subtle, though. Cyclicity of the trace is valid for ordinary matrix product, but not for the tensor product of the latter with $``$. However, the cyclicity property of Moyal product under integral (1.38) can be extended to the case of matrix-valued fields and used to prove the invariance of the action. In the limit $`\theta 0`$ one recovers the ordinary theory. It is very important to observe that while in the ordinary case $`n=1`$ corresponds to an abelian theory with field strength and transformation laws $`\delta _\lambda A_\mu =_\mu \lambda `$ (1.55) $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ (1.56) $`\delta _\lambda F_{\mu \nu }=0`$ (1.57) in the noncommutative case the commutator of two gauge transformations with parameters $`\lambda _1`$ and $`\lambda _2`$ is the gauge transformation with parameter $`\lambda _1\lambda _2\lambda _2\lambda _1`$ . This is nontrivial even in the case $`n=1`$, so this is a nonabelian theory and its features perfectly mimic the case with $`n>1`$. ##### Unitarity and causality problems As anticipated, problems arise when time-space noncommutativity is considered. The structure of Moyal product (1.32) leads to terms in the action with an infinite number of derivatives of fields. This renders a Moyal-deformed field theory non local. In particular, when time is involved in noncommutativity, nonlocality in time appears and the usual framework of quantum mechanics breaks down. In unitarity of noncommutative field theory with time-space noncommutativity has been studied. Scalar field theory deformed with time-space noncommutativity has been considered and several one loop amplitudes have been shown not to be unitary. In particular, the two point function in noncommutative $`\mathrm{\Phi }^3`$ theory has been shown not to satisfy the usual cutting rules when $`\theta ^{0i}0`$, while these rules are satisfied when only spatial noncommutativity is present. Moreover, $`22`$ scattering in noncommutative $`\mathrm{\Phi }^4`$ has been considered and again unitarity of the S-matrix is satisfied only when $`\theta ^{0i}=0`$. Recently in a series of papers a different approach to perturbative noncommutative field theories with a noncommuting time variable has been proposed. It has been argued that time-ordering is nontrivial when time is involved in noncommutativity and so a new prescription for the computation of Green functions must be given. This is different with respect to the naive Feynman rules obtained by multiplying the usual vertices by a phase factor (see (1.45)). It has been shown that in this framework unitarity is preserved when the lagrangian of the theory is hermitian. In causality of scattering processes in noncommutative field theory with time-space noncommutativity has been investigated. In particular, $`22`$ tree-level scattering amplitudes for massless scalars with $`\mathrm{\Phi }^4`$ interaction in a two-dimensional noncommutative spacetime have been computer there. The ordinary result for the $`22`$ amplitude is $$i=ig$$ (1.58) where $`g`$ is the coupling constant of the theory. In the noncommutative case with $`[t,x]=i\theta `$, because of the phases appearing in front of the vertices, one obtains instead $$ig[\mathrm{cos}(p_1p_2)\mathrm{cos}(p_3p_4)+23+24]$$ (1.59) where $`p_1,\mathrm{},p_4`$ are the two-momenta of the particles satisfying the conservation law $`_{i=1}^4p_i=0`$ (all momenta are incoming) and the wedge product is defined as $`pq\theta (p^0q^1p^1q^0)`$. In the center of mass frame (1.59) becomes $$ig[\mathrm{cos}(4p^2\theta )+2]$$ (1.60) Given an incoming gaussian wave packet $$\varphi _{in}(p)\left(e^{\frac{1}{\lambda }(pp_0)^2}+e^{\frac{1}{\lambda }(p+p_0)^2}\right)$$ (1.61) the outgoing wave packet, in the limit $`p_0\lambda ^{\frac{1}{2}}\frac{1}{p_0\theta }`$, $`\lambda \theta 1`$, is expressed as follows $$\mathrm{\Phi }_{out}(x)g\left[F(x;\theta ,\lambda ,p_0)+4\sqrt{\lambda }e^{\lambda \frac{x^2}{4}}e^{ip_0x}+F(x;\theta ,\lambda ,p_0)\right]+(p_0p_0)$$ (1.62) where $`F(x;\theta ,\lambda ,p_0)`$ represents a packet concentrated at $`x=8p_0\theta `$. So the incoming wave-packet, in the special high energy limit considered, splits into three parts, one concentrated at $`x=8p_0\theta `$, one at $`x=0`$ and one at $`x=8p_0\theta `$. All three propagate towards $`x\mathrm{}`$. In the center of mass frame scattering can be seen as bouncing on a wall. The first packet is an advanced one, which means that it leaves the wall much before the arrival of the incoming packet. The third term instead corresponds to a delayed wave, appearing well after the arrival of the incoming packet. These two terms suggest an interpretation of the noncommutative particle as a rigid rod. Both terms originate from the phase factor due to noncommutativity. The second term is not interesting, since it is neither significantly delayed or advanced. The advance by itself is not an indication of acausality. A nonrelativistic example of this is the reflection of a rigid rod of length L oriented along the direction of motion. The center of mass of the rod appears to reflect before reaching the wall, but the event is not acausal. However, there is a problem when both causality and Lorentz invariance are considered. In our case the advance increases with energy, therefore the rod seems to expand instead of Lorentz-contract at growing energies. This bizarre behavior is a sign of the inconsistency of a field theory with time-space noncommutativity. When only space-space noncommutativity is considered (in a 2+1 dimensional field theory), the effect of the phase in front of the vertices is to let outgoing scattered waves originate from the diplaced position $`y=\frac{1}{2}\theta p_x`$. This again suggests the interpretation of the incident particles as extended rods of size $`\theta p`$, but orthogonally oriented with respect to their momentum. In this case there is no violation of causality. In it has been shown that in noncommutative field theories with time-space noncommutativity tachyonic particles are produced. This gives a physical interpretation of the perturbative breakdown of unitarity. Moreover, in this paper a quantitative study of various locality and causality properties of noncommutative field theories at the quantum level has been performed. In collaboration with M.T. Grisaru, O. Lechtenfeld, L. Mazzanti, S. Penati and A. Popov I have also addressed the problem of acausality in noncommutative field theory in . We have conjectured that in a noncommutative two-dimensional field theory that is classically integrable, i.e. it has an infinite number of conserved charges, acausality may disappear. In we have shown that this is indeed the case for the noncommutative integrable sine-Gordon model, whose S-matrix is factorized, as expected for an integrable system, and causal. These results will be discussed in detail in chapter 2. The noncommutative generalization of the sine-Gordon system characterized by a well-defined S-matrix is not the natural one, though. In the next section I will begin to explore some possible noncommutative versions of the scalar theory that differ from the natural one considered in this section. In the last part of section 1.2.1 I will discuss time-space noncommutativity from the string theory point of view. There we will see that the ill-defined field theories with time-space noncommutativity do not arise as consistent limits of string theory. However, there exists a limit where one obtains a theory of open strings living in a noncommutative spacetime (NCOS). #### 1.1.3 Other possible deformations: The free scalar field theory example Up to now I have considered the natural deformation of a given field theory, obtained by simply replacing ordinary products with $``$ products everywhere in the action (and in the definitions of the field strength and gauge transformations in the Yang-Mills case). I have discussed some of the peculiar properties of noncommutative field theories obtained in such a way and noted that some of these are not welcome in a reasonable field theory. The natural deformation is not the only way to proceed. Given an ordinary field theory we can more generally define a noncommutative deformation of it as a theory written in terms of $``$ products that reproduces the original commutative theory in the limit $`\theta 0`$. Of course the natural deformation is included in this definition, but different deformations can be constructed, just by adding new terms that vanish in the limit $`\theta 0`$. I would like to discuss a very simple example, the two-dimensional free massless scalar field theory. We have previously noted that quadratic terms in the action are not modified by Moyal product. So the natural deformation of a free scalar field theory is trivial. There are more possible deformations, though, that are highly nontrivial and very interesting indeed. In ordinary geometry we can consider the element $`g`$ of a nonabelian group $`𝒢`$. With this we can construct the principal chiral model action $`S_{PC}={\displaystyle d^2x_\mu g^1^\mu g}`$ (1.63) The corresponding equation of motion is given by $`_\mu \left(g^1^\mu g\right)=0`$ (1.64) In the case of the abelian group $`𝒢=U(1)`$ we can parametrize $`g=e^{i\alpha \varphi }`$ and see that the action reduces to the free massless scalar field action. In the nonabelian case instead the model is nontrivial. It is possible to add to the principal chiral action a new term, called Wess-Zumino (WZ) term, that is written in terms of a commuting parameter $`\rho [0,1]`$. The resulting theory is called Wess-Zumino-Witten (WZW) model and its action in a $`(+,)`$ signature is $`S_{WZW}={\displaystyle \frac{1}{2}}{\displaystyle d^2x_\mu g^1^\mu g}{\displaystyle \frac{1}{3}}{\displaystyle d^2x𝑑\rho ϵ^{\mu \nu \sigma }\widehat{g}^1_\mu \widehat{g}\widehat{g}^1_\nu \widehat{g}\widehat{g}^1_\sigma \widehat{g}}`$ (1.65) (1.66) where we have introduced the homotopy path $`\widehat{g}(\rho )`$ such that $`\widehat{g}(0)=1`$, $`\widehat{g}(1)=g`$. The variation of the WZ term is a total derivative in $`\rho `$, so the equations of motion obtained from this action are truly two-dimensional and given by $`\overline{}\left(g^1g\right)=0`$ (1.67) where we have defined $`_0+_1`$ and $`\overline{}_0_1`$. For an abelian $`U(1)`$ group the WZ term vanishes and the WZW model reduces to the free massless scalar theory. We have previously observed that in noncommutative geometry the $`U(1)`$ group is no longer abelian. Parametrizing the element of noncommutative $`U(1)`$ as $`g=e_{}^{i\alpha \varphi }`$ (1.68) we can define noncommutative $`U(1)`$ principal chiral and WZW models , just substituting $``$ products everywhere in the given actions and considering $`g`$ in the noncommutative $`U(1)`$ group. Both models reduce to the free massless scalar theory in the commutative limit, so they are nontrivial noncommutative generalizations of it. Since many possible noncommutative generalizations of a single ordinary field theory can be constructed, it is natural to wonder whether one of these may be “better” than the others so that it could be chosen as “the” noncommutative version of the original theory. First of all, criteria should be given to decide whether one generalization is better than the other. There are a certain number of properties that render a field theory a “good” theory. For instance symmetries, classical integrability, causality, unitarity, renormalizability, quantum integrability, absence of anomalies, dualities are properties that, when present in the ordinary theory, we would like to preserve in its noncommutative version. So a good criteria could be to find a deformation that preserves the good properties of the ordinary theory, or some of them at least. The case of a free scalar field theory is very simple and cannot give us any hint. However, it is possible to add a potential term to the scalar theory to obtain an interacting theory enjoying nice properties. In my papers , in collaboration with M. Grisaru, O. Lechtenfeld, L. Mazzanti, S. Penati and A. Popov, I have studied the possible generalizations of a very special scalar theory, the sine-Gordon system, describing a scalar field autointeracting by an oscillating potential. As will be explained in detail in chapter 2, from our analysis of noncommutative sine-Gordon the WZW-like generalization of the kinetic term for scalar fields seems to be preferred. This agrees with bosonization considerations in . #### 1.1.4 The Kontsevich product In section 1.1.1 we have seen that, in a Poisson manifold with a derivative without torsion and curvature, such that $`P=0`$,the only associative product is Moyal $``$ (1.32). If one of the three assumptions is not verified, then this product is no longer associative. Moreover, we have noticed that in flat spacetime with ordinary derivatives the assumption of constant $`P`$ is a restriction needed for the noncommutative deformation of a translation invariant theory to preserve this property. In this section I would like to consider more general situations, where spacetime may be curved and translations may not be a symmetry anymore. In particular, as I will show in detail in section 1.2.1, noncommutative geometry naturally emerges in the context of string theory, that is naturally embedded in curved backgrounds. Therefore, spacetimes with torsion and curvature should be considered. Different symmetries underlying the theory may require Poisson structures $`P`$ with a particular dependence on the coordinates. For this reasons it is interesting to consider the case when the Poisson structure $`P`$ is not constant with respect to a certain set of derivatives. One may think about relaxing the other two constraints regarding the set of derivatives chosen. Actually, at least in some cases it is possible to rewrite a Poisson structure with nonflat derivatives and a covariantly constant Poisson tensor in terms of a nonconstant Poisson tensor and flat derivatives. The superspace case we will study in the next section is an example of this. M. Kontsevich in generalized the results of Bayen et al. to the case where derivatives $``$ are without torsion and curvature but the Poisson structure is not covariantly constant $`P0`$. In this case Moyal product (1.32) is not associative anymore. However, it is possible to modify Moyal product order by order in the deformation parameter $`\mathrm{}`$ to obtain associativity. Let $`\mathrm{\Omega }`$ be a Poisson manifold, with coordinates $`\{x^\mu \}`$ and flat derivatives $`\{_\mu \}`$. Let $`P^{\mu \nu }=P^{\mu \nu }(x)`$ be the Poisson structure of this manifold, written in terms of the chosen coordinates as follows $$[x^\mu ,x^\nu ]=P^{\mu \nu }(x)$$ (1.69) First of all we recall that the definition of a Poisson structure requires associativity $$P^{\mu \rho }_\rho P^{\nu \sigma }+P^{\nu \rho }_\rho P^{\sigma \mu }+P^{\sigma \rho }_\rho P^{\mu \nu }=0$$ (1.70) that is completely equivalent to the validity of Jacobi identity for the coordinate algebra $$[[x^\mu ,x^\nu ],x^\sigma ]+[[x^\sigma ,x^\mu ],x^\nu ]+[[x^\nu ,x^\sigma ],x^\mu ]=0$$ (1.71) When $`P^{\mu \nu }`$ is constant, this requirement is trivially satisfied and does not give any further constraint on the matrix $`P`$. When a general coordinate dependence is allowed, associativity constrains the functional dependence of the matrix. When $`P^{\mu \nu }`$ is invertible, i.e. it exists $`P_{\mu \nu }^1`$ satisfying $`P^{\mu \nu }P_{\nu \rho }^1=\delta _\rho ^\mu `$, (1.70) can easily be rewritten in terms of the vanishing of the three-form $`H`$ $$H_{\mu \nu \rho }=_\mu P_{\nu \rho }+\mathrm{cycl}.=0$$ (1.72) Let us now consider Moyal product (1.32), expanded up to second order in the parameter $`\mathrm{}`$ $`fg=fg+\mathrm{}P^{\mu \nu }(x)_\mu f_\nu g+\mathrm{}^2P^{\mu \nu }(x)P^{\rho \sigma }(x)_\mu _\rho f_\nu _\sigma g+𝒪(\mathrm{}^3)`$ (1.73) (1.74) We then evaluate the quantity $`(fg)hf(gh)`$ up to second order in $`\mathrm{}`$ (note that associativity is trivially satisfied at first order): $$(fg)hf(gh)=\mathrm{}^2\left(P^{\sigma \rho }_\rho P^{\mu \nu }+P^{\mu \rho }_\rho P^{\nu \sigma }\right)_\mu f_\nu g_\sigma h+𝒪(\mathrm{}^3)$$ (1.75) So nonvanishing terms arise, because of the $`x`$ dependence of the Poisson structure. Kontsevich observed that once the trilinear dependence on the functions $`f`$, $`g`$, $`h`$ is factorized, one obtains terms with an identical structure with respect to the ones emerging in the associativity equation for $`P`$ (1.70). If one could modify the definition of the product $``$ (1.74) in such a way to obtain exactly the quantity that is constrained to be zero in (1.70), one would obtain a product associative up to order $`\mathrm{}^2`$. We have to add new terms of order $`\mathrm{}^2`$, since associativity is trivially satisfied at first order. Therefore let us define a new product $``$ by adding to $``$ a new term of order $`\mathrm{}^2`$ as follows $`fg`$ $`fg+\mathrm{}P^{\mu \nu }(x)_\mu f_\nu g+\mathrm{}^2P^{\mu \nu }(x)P^{\rho \sigma }(x)_\mu _\rho f_\nu _\sigma g`$ (1.78) $`+A\mathrm{}^2P^{\mu \rho }_\rho P^{\nu \sigma }\left(_\mu _\nu f_\sigma g_\nu f_\mu _\sigma g\right)+𝒪(\mathrm{}^3)`$ where $`A`$ is a coefficient to be determined in such a way that $`(fg)hf(gh)`$ is proportional to the constraint (1.70). One obtains $$A=\frac{1}{3}$$ (1.79) The product $``$ that we have defined up to second order in $`\mathrm{}`$ is associative if and only if the Jacobi identity for the coordinate algebra (1.71) is satisfied. Kontsevich showed that order by order in $`\mathrm{}`$ it is always possible to modify Moyal product $``$ in order to make extra terms vanish when the Jacobi identity for the coordinates is valid. So Kontsevich $``$ product is uniquely defined at any order in the deformation parameter $`\mathrm{}`$ by the requirement of associativity. #### 1.1.5 Deforming superspace Up to now, we have only considered noncommutativity of bosonic coordinates, in the form $$[x^\mu ,x^\nu ]=i\mathrm{\Theta }^{\mu \nu }(x)$$ (1.80) where $`\mathrm{\Theta }^{\mu \nu }(x)`$ is antisymmetric. Since, as it will be explained in section 2, noncommutative field theories emerge naturally in the context of string theory and string theory is only consistent in the presence of supersymmetry, it is natural to consider the problem of deforming a supersymmetric theory. It is well-known that the natural setting for discussing supersymmetric theories is a nontrivial extension of bosonic space, known as superspace, where bosonic coordinates $`x`$ are accompanied by fermionic ones, that I will generally denote with $`\theta `$. So it seems natural and compelling to ask what happens if we deform also the anticommutators between fermionic coordinates of superspace. Exactly as in the bosonic case discussed before, we would like the deformation to preserve the symmetries of our original theory, described by the group of supertranslations. Moreover, we would like the deformed algebra to be associative. In the bosonic case we have seen that this two properties in flat space were identifying Moyal product. In collaboration with D. Klemm and S. Penati, I have addressed the problem of deforming superspace in . First steps had been taken in this direction before the appearance of this paper. Nonvanishing anticommutators of fermionic coordinates have been considered in in the context of a possible fermionic substructure of spacetime. In , quantum deformations of the Poincaré supergroup were considered. In a modern noncommutative geometry context, trivial superspace deformation of supersymmetric field theories have been analysed, where only the bosonic sector of the coordinate algebra is modified . In , a Moyal-like deformation of $`d=4`$ $`N=1`$ superspace was proposed involving fermionic coordinates, that is associative and covariant with respect to supersymmetry, but does not preserve the complex conjugation rules that characterize Majorana-Weyl spinors in four dimensions. In this paper it was also shown that in general the set of chiral superfields is not closed under star products that involve fermionic coordinates. I will discuss the results in in detail in the second part of this section. In we were mainly concerned with the conditions imposed on the possible deformations of superspace by requirements such as covariance under classical translations and supertranslations, Jacobi identities, associativity of the star product and closure of the set of chiral superfields under the star product, but we wanted superspace conjugation relations to be still valid in the deformation. The motivation for this requirement relies in the fact that in a theory with $`N`$ supersymmetries formulated in superspace, the number of fermionic degrees of freedom is $`N`$ times the bosonic one. As I will show in detail, by relaxing spinor conjugation relations in the deformation, one in fact does not preserve the number of supersymmetries. The main results in are that it is possible to fulfill all the requirements in a $`d=4`$ $`N=1`$ Minkowski superspace, even if the contraints imposed on the supercoordinate algebra are strong and only $`[x,\theta ]`$ and $`[x,x]`$ can be turned on. Moreover, in the same paper we have shown that euclidean signature is less restrictive and allows for a nonanticommutative superspace with $`\{\theta ,\theta \}`$ different from zero. The results obtained in my paper will be discussed in section 1.1.6. Since then a lot of progress has been done in understanding superspace deformations. Non(anti)commutative superspaces have been shown to emerge in a superstring theory context, in the presence of Ramond-Ramond (RR) backgrounds . I will discuss these results in section 1.2.4. Furthermore, in the paper by N. Seiberg , a deformed superspace that only preserves $`N=\frac{1}{2}`$ supersymmetric of the original $`N=1`$ has been introduced. I will review the properties of this deformation in section 1.1.7 and compare with the ones obtained in my paper. Many deformations of superspace field theories have been studied and their quantum properties have been discussed. I will give a brief summary of the main results obtained in the second part of section 1.1.7. I will not give an introduction to supersymmetry and its superspace formulation. For an introduction to this topics, I suggest the books and the review paper . ##### The Ferrara-Lledò proposal In the authors consider the problem of generalizing Moyal product (1.32) to $`d=4`$ $`N=1`$ superspace. This is described by the set of superspace coordinates $`Z^A=(x^{\alpha \dot{\alpha }},\theta ^\alpha ,\overline{\theta }^{\dot{\alpha }})`$, where $`x^{\alpha \dot{\alpha }}`$ are four real bosonic coordinates and $`\theta ^\alpha `$, $`\overline{\theta }^{\dot{\alpha }}`$ are two–component complex Weyl fermions. The conjugation rule $`\overline{\theta }^{\dot{\alpha }}=(\theta ^\alpha )^{}`$ follows from the requirement to have a four component Majorana fermion (we use conventions of Superspace ). In the standard (anti)commutative superspace the algebra of the coordinates is $`\{\theta ^\alpha ,\theta ^\beta \}=\{\overline{\theta }^{\dot{\alpha }},\overline{\theta }^{\dot{\beta }}\}=\{\theta ^\alpha ,\overline{\theta }^{\dot{\alpha }}\}=0`$ $`[x^{\alpha \dot{\alpha }},\theta ^\beta ]=[x^{\alpha \dot{\alpha }},\overline{\theta }^{\dot{\beta }}]=0`$ $`[x^{\alpha \dot{\alpha }},x^{\beta \dot{\beta }}]=0`$ (1.81) and it is trivially covariant under the superPoincaré group. The subgroup of the classical (super)translations (spacetime translations and supersymmetry transformations) $`\theta ^\alpha =\theta ^\alpha +ϵ^\alpha `$ $`\overline{\theta }^{\dot{\alpha }}=\overline{\theta }^{\dot{\alpha }}+\overline{ϵ}^{\dot{\alpha }}`$ $`x^{\alpha \dot{\alpha }}=x^{\alpha \dot{\alpha }}+a^{\alpha \dot{\alpha }}{\displaystyle \frac{i}{2}}\left(ϵ^\alpha \overline{\theta }^{\dot{\alpha }}+\overline{ϵ}^{\dot{\alpha }}\theta ^\alpha \right)`$ (1.82) is generated by two complex charges $`Q_\alpha `$ ($`\overline{Q}_{\dot{\alpha }}=Q_\alpha ^{}`$) and the four–momentum $`P_{\alpha \dot{\alpha }}`$ subjected to $$\{Q_\alpha ,Q_\beta \}=\{\overline{Q}_{\dot{\alpha }},\overline{Q}_{\dot{\beta }}\}=0,\{Q_\alpha ,\overline{Q}_{\dot{\alpha }}\}=P_{\alpha \dot{\alpha }}$$ (1.83) Representations of supersymmetry are given by superfields $`V(x^{\alpha \dot{\alpha }},\theta ^\alpha ,\overline{\theta }^{\dot{\alpha }})`$ whose components are obtained by expanding $`V`$ in powers of the spinorial coordinates. The set of superfields is closed under the standard product of functions. The product of two superfields is (anti)commutative, $`VW=(1)^{deg(V)deg(W)}WV`$, and associative, $`(KV)W=K(VW)`$. The set of superspace covariant derivatives is given by $`_A=(_{\alpha \dot{\alpha }},D_\alpha ,\overline{D}_{\dot{\alpha }})`$. In both this section and the following one (so in the papers and ) the nonchiral representation of supersymmetry has been chosen, where $`D_\alpha =_\alpha +{\displaystyle \frac{i}{2}}\overline{\theta }^{\dot{\alpha }}_{\alpha \dot{\alpha }};\overline{D}_{\dot{\alpha }}=\overline{}_{\dot{\alpha }}+{\displaystyle \frac{i}{2}}\theta ^\alpha _{\alpha \dot{\alpha }}`$ (1.84) Superspace geometry is nontrivial, because of the presence of a nonvanishing torsion $$\{D_\alpha ,\overline{D}_{\dot{\alpha }}\}=i_{\alpha \dot{\alpha }}$$ (1.85) To extend the construction of Moyal product to superspace, one has first to introduce a superspace Poisson structure generalizing (1.20). The authors propose $$\{\mathrm{\Phi },\mathrm{\Psi }\}=\mathrm{\Phi }\stackrel{}{}_{\alpha \dot{\alpha }}P^{\alpha \dot{\alpha }\beta \dot{\beta }}\stackrel{}{}_{\beta \dot{\beta }}\mathrm{\Psi }+\mathrm{\Phi }\stackrel{}{D}_\alpha P^{\alpha \beta }\stackrel{}{D}_\beta \mathrm{\Psi }$$ (1.86) where $`P^{\alpha \dot{\alpha }\beta \dot{\beta }}`$ and $`P^{\alpha \beta }`$ are constant matrices. This Poisson structure is manifestly covariant with respect to supersymmetry. Moreover, it is associative since it involves only $`_{\alpha \dot{\alpha }}`$ and $`D_\alpha `$ and not $`\overline{D}_{\dot{\alpha }}`$. Associativity would be broken if the whole set of superspace covariant derivatives had appeared, because of the nontrivial superspace torsion (1.85). With this Poisson structure the authors construct a Moyal-like product in superspace as follows $$\mathrm{\Phi }\mathrm{\Psi }=\mathrm{\Phi }\mathrm{exp}\left(\mathrm{}\left(\stackrel{}{}_{\alpha \dot{\alpha }}P^{\alpha \dot{\alpha }\beta \dot{\beta }}\stackrel{}{}_{\beta \dot{\beta }}+\stackrel{}{D}_\alpha P^{\alpha \beta }\stackrel{}{D}_\beta \right)\right)\mathrm{\Psi }$$ (1.87) This is an associative product, as one can easily deduce by extending to superspace the discussion in section 1.1.1. If we consider the special case when the superfields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are identified with the supercoordinates themselves, we obtain the following anticommutation relations $`\{\theta ^\alpha ,\theta ^\beta \}=P^{\alpha \beta }`$ (1.88) $`\{\overline{\theta }^{\dot{\alpha }},\overline{\theta }^{\dot{\beta }}\}=0`$ (1.89) They are not consistent with the $`d=4`$ $`N=1`$ superspace conjugation relation $$\left(\theta ^\alpha \right)^{}=\overline{\theta }^{\dot{\alpha }}$$ (1.90) Even if in principle a noncommutative deformation is not required to respect the complex structure present on the original space, the relation (1.90) is needed in the original space for the fermionic degrees of freedom to be equal to the bosonic ones. By relaxing it in the deformation, one modifies (doubles) the number of supersymmetries. In the paper it was also observed that the class of chiral superfields, defined by the relation $`\overline{D}_{\dot{\alpha }}\mathrm{\Phi }=0`$, is not closed under the product (1.87). Already at first order in the $`\mathrm{}`$ expansion nontrivial terms arise in $`\overline{D}_{\dot{\alpha }}\left(\mathrm{\Phi }\mathrm{\Psi }\right)`$, where $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are both chiral, because of the nontrivial superspace torsion (1.85). Since the simplest superspace field theories are written in terms of chiral superfields, the lack of closure for the chiral class is a serious obstruction in constructing deformations of known supersymmetric theories. #### 1.1.6 Non(anti)commutative superspace In this section I will discuss the results that I obtained in , in collaboration with D. Klemm and S. Penati, regarding supersymmetric associative deformations of $`d=4`$ $`N=1`$ Minkowski and $`N=2`$ euclidean superspace. ##### Supersymmetric deformations of $`N=1`$ $`d=4`$ superpace In , a more systematical approach with respect to has been followed to determine the most general non(anti)commutative geometry in $`N=1`$ four dimensional superspace, invariant under the classical supertranslation group and associative. As I have anticipated before, the deformation will be required to preserve the complex conjugation relations that are valid in ordinary superspace. We will consider the supercoordinates $`Z^A`$, generically satisfying the non(anti)commutative algebra $$[Z^A,Z^B\}=P^{AB}(Z)$$ (1.91) where we have introduced the (anti)commutator $$[F_A,G_B\}F_AG_B()^{ab}G_BF_A$$ (1.92) that is a commutator if at least one of the two indices $`A`$, $`B`$ is a vector and an anticommutator otherwise. Let us consider the transformation $`ZZ^{}`$. This is the generic symmetry of the ordinary theory that we would like to preserve in the deformation. Exactly as in the bosonic case discussed in section 1.1.2, we will require that the functional dependence of the non(anti)commutative algebra is not modified under the transformation $$[Z_{}^{}{}_{}{}^{A},Z_{}^{}{}_{}{}^{B}\}=P^{AB}(Z^{})$$ (1.93) Since, as discussed in section 1.1.1, bosonic noncommutative deformations break “particle” Lorentz invariance, we will not worry about this symmetry in our superspace generalization. We will only take into consideration the supertranslation group, containing ordinary bosonic translations and supersymmetry transformations. For $`N=1`$ $`d=4`$ superspace conventions, we refer to (see summary in the previous section). In order to define a non(anti)commutative superspace, we consider the most general structure of the algebra for a set of four bosonic real coordinates and a complex two–component Weyl spinor with $`(\theta ^\alpha )^{}=\overline{\theta }^{\dot{\alpha }}`$ $`\{\theta ^\alpha ,\theta ^\beta \}=𝒜^{\alpha \beta }(x,\theta ,\overline{\theta }),\{\overline{\theta }^{\dot{\alpha }},\overline{\theta }^{\dot{\beta }}\}=\overline{𝒜}^{\dot{\alpha }\dot{\beta }}(x,\theta ,\overline{\theta })`$ $`\{\theta ^\alpha ,\overline{\theta }^{\dot{\alpha }}\}=^{\alpha \dot{\alpha }}(x,\theta ,\overline{\theta })`$ $`[x^{\underset{¯}{a}},\theta ^\beta ]=i𝒞^{\underset{¯}{a}\beta }(x,\theta ,\overline{\theta }),[x^{\underset{¯}{a}},\overline{\theta }^{\dot{\beta }}]=i\overline{𝒞}^{\underset{¯}{a}\dot{\beta }}(x,\theta ,\overline{\theta })`$ $`[x^{\underset{¯}{a}},x^{\underset{¯}{b}}]=i𝒟^{\underset{¯}{a}\underset{¯}{b}}(x,\theta ,\overline{\theta })`$ (1.94) Here, $`𝒜,,𝒞,𝒟`$ are local functions of the superspace variables and we have defined $`\overline{𝒜}^{\dot{\alpha }\dot{\beta }}(𝒜^{\alpha \beta })^{}`$, $`\overline{𝒞}^{\underset{¯}{a}\dot{\beta }}(𝒞^{\underset{¯}{a}\beta })^{}`$. From the conjugation rules for the coordinates it follows also $`\left(^{\alpha \dot{\alpha }}\right)^{}=^{\alpha \dot{\alpha }}`$ and $`\left(𝒟^{\underset{¯}{a}\underset{¯}{b}}\right)^{}=𝒟^{\underset{¯}{a}\underset{¯}{b}}`$. To implement (1.94) to be the algebra of the coordinates of a non(anti)commutative $`N=1`$ superspace we require its invariance under the group of space translations and supertranslations (1.82). As before, we restrict our analysis to the case of an undeformed group where the parameters $`a^{\alpha \dot{\alpha }}`$, $`ϵ^\alpha `$ and $`\overline{ϵ}^{\dot{\alpha }}`$ in (1.82) are kept (anti)commuting <sup>2</sup><sup>2</sup>2More general constructions of non(anti)commutative geometries in grassmannian spaces have been considered, where also the algebra of the parameters is deformed .. As in the bosonic case discussed in section 1.1.1, we require the functional dependence of the $`𝒜,,𝒞,𝒟`$ in (1.94) to be the same at any point of the supermanifold. To work out explicitly the constraints which follow, we perform a (super)translation (1.82) on the coordinates and compute the algebra of the new coordinates in terms of the old ones. We find that the functions appearing in (1.94) are constrained by the following set of independent equations $$𝒜^{\alpha \beta }(x^{},\theta ^{},\overline{\theta }^{})=𝒜^{\alpha \beta }(x,\theta ,\overline{\theta }),^{\alpha \dot{\alpha }}(x^{},\theta ^{},\overline{\theta }^{})=^{\alpha \dot{\alpha }}(x,\theta ,\overline{\theta })$$ (1.95) $$𝒞^{\alpha \dot{\alpha }\beta }(x^{},\theta ^{},\overline{\theta }^{})=𝒞^{\alpha \dot{\alpha }\beta }(x,\theta ,\overline{\theta })\frac{1}{2}ϵ^\alpha ^{\beta \dot{\alpha }}(x,\theta ,\overline{\theta })\frac{1}{2}\overline{ϵ}^{\dot{\alpha }}𝒜^{\alpha \beta }(x,\theta ,\overline{\theta })$$ (1.96) $`𝒟^{\alpha \dot{\alpha }\beta \dot{\beta }}(x^{},\theta ^{},\overline{\theta }^{})=𝒟^{\alpha \dot{\alpha }\beta \dot{\beta }}(x,\theta ,\overline{\theta })`$ $`{\displaystyle \frac{i}{2}}\left(ϵ^\beta \overline{𝒞}^{\alpha \dot{\alpha }\dot{\beta }}(x,\theta ,\overline{\theta })+\overline{ϵ}^{\dot{\beta }}𝒞^{\alpha \dot{\alpha }\beta }(x,\theta ,\overline{\theta })ϵ^\alpha \overline{𝒞}^{\beta \dot{\beta }\dot{\alpha }}(x,\theta ,\overline{\theta })\overline{ϵ}^{\dot{\alpha }}𝒞^{\beta \dot{\beta }\alpha }(x,\theta ,\overline{\theta })\right)`$ $`{\displaystyle \frac{i}{4}}(ϵ^\alpha \overline{𝒜}^{\dot{\alpha }\dot{\beta }}(x,\theta ,\overline{\theta })ϵ^\beta +ϵ^\alpha ^{\beta \dot{\alpha }}(x,\theta ,\overline{\theta })\overline{ϵ}^{\dot{\beta }}`$ $`+\overline{ϵ}^{\dot{\alpha }}^{\alpha \dot{\beta }}(x,\theta ,\overline{\theta })ϵ^\beta +\overline{ϵ}^{\dot{\alpha }}𝒜^{\alpha \beta }(x,\theta ,\overline{\theta })\overline{ϵ}^{\dot{\beta }})`$ (1.97) together with their hermitian conjugates. Looking for the most general local solution brings us to the following algebra for a non(anti)commutative geometry in Minkowski superspace consistent with (super)translations $`\{\theta ^\alpha ,\theta ^\beta \}`$ $`=A^{\alpha \beta },\{\overline{\theta }^{\dot{\alpha }},\overline{\theta }^{\dot{\beta }}\}=\overline{A}^{\dot{\alpha }\dot{\beta }},\{\theta ^\alpha ,\overline{\theta }^{\dot{\alpha }}\}=B^{\alpha \dot{\alpha }}`$ $`[x^{\alpha \dot{\alpha }},\theta ^\beta ]`$ $`=i𝒞^{\alpha \dot{\alpha }\beta }(\theta ,\overline{\theta })`$ $`[x^{\alpha \dot{\alpha }},\overline{\theta }^{\dot{\beta }}]`$ $`=i\overline{𝒞}^{\alpha \dot{\alpha }\dot{\beta }}(\theta ,\overline{\theta })`$ $`[x^{\alpha \dot{\alpha }},x^{\beta \dot{\beta }}]`$ $`=i𝒟^{\alpha \dot{\alpha }\beta \dot{\beta }}(\theta ,\overline{\theta })`$ (1.98) where $`𝒞^{\alpha \dot{\alpha }\beta }(\theta ,\overline{\theta })=C^{\alpha \dot{\alpha }\beta }{\displaystyle \frac{1}{2}}\theta ^\alpha B^{\beta \dot{\alpha }}{\displaystyle \frac{1}{2}}\overline{\theta }^{\dot{\alpha }}A^{\alpha \beta }`$ (1.99) $`𝒟^{\alpha \dot{\alpha }\beta \dot{\beta }}(\theta ,\overline{\theta })=D^{\alpha \dot{\alpha }\beta \dot{\beta }}{\displaystyle \frac{i}{2}}\left(\theta ^\beta \overline{C}^{\alpha \dot{\alpha }\dot{\beta }}\overline{\theta }^{\dot{\alpha }}C^{\beta \dot{\beta }\alpha }\theta ^\alpha \overline{C}^{\beta \dot{\beta }\dot{\alpha }}+\overline{\theta }^{\dot{\beta }}C^{\alpha \dot{\alpha }\beta }\right)`$ $`{\displaystyle \frac{i}{4}}\left(\theta ^\alpha \overline{A}^{\dot{\alpha }\dot{\beta }}\theta ^\beta +\theta ^\alpha B^{\beta \dot{\alpha }}\overline{\theta }^{\dot{\beta }}+\overline{\theta }^{\dot{\alpha }}B^{\alpha \dot{\beta }}\theta ^\beta +\overline{\theta }^{\dot{\alpha }}A^{\alpha \beta }\overline{\theta }^{\dot{\beta }}\right)`$ (1.100) and $`A`$, $`B`$, $`C`$ and $`D`$ are constant functions. We note that, while invariance under spacetime translations necessarily requires the non(anti)commutation functions to be independent of the $`x`$ coordinates, as we have seen in section 1.1.1, invariance under supersymmetry is less restrictive and allows for a particular dependence on the spinorial coordinates. On the algebra of smooth functions of superspace variables we can formally define a graded bracket which reproduces the fundamental algebra (1.98) when applied to the coordinates. In the case of bosonic Minkowski spacetime, the noncommutative algebra (1.80) can be obtained by interpreting the l.h.s. of this relation as the Poisson bracket of classical commuting variables, where, for generic functions of spacetime, the Poisson bracket is defined as $`\{f,g\}_P=i\mathrm{\Theta }^{\mu \nu }_\mu f_\nu g`$. Generalizing to Minkowski superspace, the graded bracket must be constructed as a bidifferential operator with respect to the superspace variables. Using covariant derivatives $`D_A(D_\alpha ,\overline{D}_{\dot{\alpha }},_{\alpha \dot{\alpha }})`$, for generic functions $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ of the superspace coordinates we define the bidifferential operator $$\{\mathrm{\Phi },\mathrm{\Psi }\}_P=\mathrm{\Phi }\stackrel{}{D}_AP^{AB}\stackrel{}{D}_B\mathrm{\Psi }$$ (1.101) where $$P^{AB}\left(\begin{array}{ccc}P^{\alpha \beta }& P^{\alpha \dot{\beta }}& P^{\alpha \underset{¯}{b}}\\ P^{\dot{\alpha }\beta }& P^{\dot{\alpha }\dot{\beta }}& P^{\dot{\alpha }\underset{¯}{b}}\\ P^{\underset{¯}{a}\beta }& P^{\underset{¯}{a}\dot{\beta }}& P^{\underset{¯}{a}\underset{¯}{b}}\end{array}\right)=\left(\begin{array}{ccc}A^{\alpha \beta }& B^{\alpha \dot{\beta }}& iC^{\beta \dot{\beta }\alpha }\\ B^{\dot{\alpha }\beta }& \overline{A}^{\dot{\alpha }\dot{\beta }}& i\overline{C}^{\beta \dot{\beta }\dot{\alpha }}\\ iC^{\alpha \dot{\alpha }\beta }& i\overline{C}^{\alpha \dot{\alpha }\dot{\beta }}& iD^{\alpha \dot{\alpha }\beta \dot{\beta }}\end{array}\right)$$ (1.102) is a constant graded symplectic supermatrix satisfying $`P^{BA}=(1)^{(a+1)(b+1)}P^{AB}`$, $`a`$ denoting the grading of $`A`$. It is easy to verify that applying this operator to the superspace coordinates we obtain (1.98). Alternatively, one can express the graded brackets (1.101) in terms of torsion free, noncovariant spinorial derivatives $`_A(_\alpha ,\overline{}_{\dot{\alpha }},_{\alpha \dot{\alpha }})`$ so obtaining a matrix $`\stackrel{~}{P}^{AB}`$ explicitly dependent on $`(\theta ,\overline{\theta })`$. The bracket (1.101) is rewritten as follows $$\{\mathrm{\Phi },\mathrm{\Psi }\}_P=\mathrm{\Phi }\stackrel{}{}_A\stackrel{~}{P}^{AB}\stackrel{}{}_B\mathrm{\Psi }$$ (1.103) where $$\stackrel{~}{P}^{AB}(\theta ,\overline{\theta })\left(\begin{array}{ccc}A^{\alpha \beta }& B^{\alpha \dot{\beta }}& i𝒞^{\beta \dot{\beta }\alpha }(\theta ,\overline{\theta })\\ B^{\dot{\alpha }\beta }& \overline{A}^{\dot{\alpha }\dot{\beta }}& i\overline{𝒞}^{\beta \dot{\beta }\dot{\alpha }}(\theta ,\overline{\theta })\\ i𝒞^{\alpha \dot{\alpha }\beta }(\theta ,\overline{\theta })& i\overline{𝒞}^{\alpha \dot{\alpha }\dot{\beta }}(\theta ,\overline{\theta })& i𝒟^{\alpha \dot{\alpha }\beta \dot{\beta }}(\theta ,\overline{\theta })\end{array}\right)$$ (1.104) and the functions $`𝒞`$, $`\overline{𝒞}`$ and $`𝒟`$ are given in (1.100). The latter formulation of the graded bracket is not manifestly covariant, but it is however very useful, since it makes clear that Kontsevich procedure outlined in section 1.1.4 can be generalized to superspace to construct an associative deformation of the product between superfields. All the results from now on can be written in both ways. In general the covariant formulation is preferred because it naturally leads to a geometrical interpretation of the results. ##### Deformation of the supersymmetry algebra It is important to note that the non(anti)commutative extension given in (1.98) in general deforms the supersymmetry algebra. In the standard case, defining $`Q_A(Q_\alpha ,\overline{Q}_{\dot{\alpha }},i_{\alpha \dot{\alpha }})`$, the supersymmetry algebra can be written as $`[Q_A,Q_B\}=iT_{AB}^{}{}_{}{}^{C}Q_C,[D_A,D_B\}=T_{AB}^{}{}_{}{}^{C}D_C`$ $`[Q_A,D_B\}=0`$ (1.105) where $`T_{AB}^{}{}_{}{}^{C}`$ is the torsion of the flat superspace ($`T_{\alpha \dot{\beta }}^{}{}_{}{}^{\underset{¯}{c}}=T_{\dot{\beta }\alpha }^{}{}_{}{}^{\underset{¯}{c}}=i\delta _\alpha ^\gamma \delta _{\dot{\beta }}^{\dot{\gamma }}`$ are the only nonzero components). Turning on non(anti)commutativity in superspace leads instead to $`[Q_A,Q_B\}=iT_{AB}^{}{}_{}{}^{C}Q_C+R_{AB}^{}{}_{}{}^{CD}Q_CQ_D`$ $`[D_A,D_B\}=T_{AB}^{}{}_{}{}^{C}D_C+R_{AB}^{}{}_{}{}^{CD}D_CD_D`$ $`[Q_A,D_B\}=R_{AB}^{}{}_{}{}^{CD}Q_CD_D`$ (1.106) where $`T_{AB}^{}{}_{}{}^{C}`$ is still the torsion of the flat superspace, while $$R_{AB}^{}{}_{}{}^{CD}=\frac{1}{8}P^{MN}T_{M[A}^{}{}_{}{}^{C}T_{B)N}^{}{}_{}{}^{D}$$ (1.107) ($`[ab)`$ means antisymmetrization when at least one of the indices is a vector index, symmetrization otherwise) is a curvature tensor whose presence is a direct consequence of the non(anti)commutation of the grassmannian coordinates. Its nonvanishing components are $`R_{\alpha \beta }^{}{}_{}{}^{\underset{¯}{c}\underset{¯}{d}}={\displaystyle \frac{1}{8}}P^{\dot{\gamma }\dot{\delta }}\delta _{(\alpha }^\gamma \delta _{\beta )}^\delta ,R_{\dot{\alpha }\dot{\beta }}^{}{}_{}{}^{\underset{¯}{c}\underset{¯}{d}}={\displaystyle \frac{1}{8}}P^{\gamma \delta }\delta _{(\dot{\alpha }}^{\dot{\gamma }}\delta _{\dot{\beta })}^{\dot{\delta }}`$ $`R_{\alpha \dot{\beta }}^{}{}_{}{}^{\underset{¯}{c}\underset{¯}{d}}=R_{\dot{\beta }\alpha }^{}{}_{}{}^{\underset{¯}{c}\underset{¯}{d}}={\displaystyle \frac{1}{8}}\left(P^{\gamma \dot{\delta }}\delta _\alpha ^\delta \delta _{\dot{\beta }}^{\dot{\gamma }}+P^{\delta \dot{\gamma }}\delta _\alpha ^\gamma \delta _{\dot{\beta }}^{\dot{\delta }}\right)`$ (1.108) I would like to stress that the curvature terms deforming the supersymmetry algebra are quadratic in bosonic derivatives and have no effect on the supercoordinate-algebra. Therefore their presence is not in disagreement with consistency of the supercoordinate algebra with supersymmetry. Since the terms proportional to the curvature in (1.106) are quadratic in supersymmetry charges and covariant derivatives, we can define new graded brackets $`[Q_A,Q_B\}_qQ_AQ_B(1)^{ab}[\delta _{B}^{}{}_{}{}^{C}\delta _{A}^{}{}_{}{}^{D}+(1)^{ab}R_{AB}^{}{}_{}{}^{CD}]Q_CQ_D`$ (1.109) (1.110) and analogous ones for $`[D_A,D_B\}_q`$ and $`[Q_A,D_B\}_q`$, which satisfy the standard algebra (1.105). The new brackets can be interpreted as a quantum deformation associated to a $`q`$–parameter which in this case is a rank–four tensor $$q_{AB}^{}{}_{}{}^{CD}\delta _{B}^{}{}_{}{}^{C}\delta _{A}^{}{}_{}{}^{D}+(1)^{ab}R_{AB}^{}{}_{}{}^{CD}$$ (1.111) ##### Associativity and the geometry of deformed superspace Given the bidifferential operator (1.101) associated to the noncommutative supergeometry defined in (1.98) it is easy to prove the following identities $`\{\mathrm{\Phi },\mathrm{\Psi }\}_P=(1)^{1+deg(\mathrm{\Phi })deg(\mathrm{\Psi })}\{\mathrm{\Psi },\mathrm{\Phi }\}_P`$ $`\{c\mathrm{\Phi },\mathrm{\Psi }\}_P=c\{\mathrm{\Phi },\mathrm{\Psi }\}_P,\{\mathrm{\Phi },c\mathrm{\Psi }\}_P=(1)^{deg(c)deg(\mathrm{\Phi })}c\{\mathrm{\Phi },\mathrm{\Psi }\}_P`$ $`\{\mathrm{\Phi }+\mathrm{\Psi },\mathrm{\Omega }\}_P=\{\mathrm{\Phi },\mathrm{\Omega }\}_P+\{\mathrm{\Psi },\mathrm{\Omega }\}_P`$ (1.112) The operator $`\{,\}_P`$ will then be promoted to a graded Poisson structure on superspace if and only if Jacobi identities hold $`\{\mathrm{\Phi },\{\mathrm{\Psi },\mathrm{\Omega }\}_P\}_P+(1)^{deg(\mathrm{\Phi })[deg(\mathrm{\Psi })+deg(\mathrm{\Omega })]}\{\mathrm{\Psi },\{\mathrm{\Phi },\mathrm{\Omega }\}_P\}_P`$ $`+(1)^{deg(\mathrm{\Omega })[deg(\mathrm{\Phi })+deg(\mathrm{\Psi })]}\{\mathrm{\Omega },\{\mathrm{\Phi },\mathrm{\Psi }\}_P\}_P=0`$ (1.113) for any triplet of functions of superspace variables. This property is equivalent to associativity of the fundamental algebra (1.98). Since the latter is nontrivial (coordinate-dependent commutators appear), (1.113) is not in general satisfied. Indeed, imposing (1.113) yields the nontrivial conditions $`P^{AR}P^{BS}T_{SR}^C(1)^{c+b(c+a+r)}+P^{BR}P^{CS}T_{SR}^A(1)^{a+c(a+b+r)}`$ $`+P^{CR}P^{AS}T_{SR}^B(1)^{b+a(b+c+r)}=0`$ (1.114) $$(1)^{bm}P^{AM}P^{BN}R_{MN}^{}{}_{}{}^{CD}=0$$ (1.115) where the torsion $`T_{AB}^{}{}_{}{}^{C}`$ and the curvature $`R_{AB}^{}{}_{}{}^{CD}`$ have been introduced in (1.105) and (1.106). Equation (1.114) is the covariant superspace generalization of bosonic associativity constraint (1.70) (the analogy with (1.70) is clearer when it is rewritten in terms of the “noncovariant” Poisson structure $`\stackrel{~}{P}`$ given in (1.103)). As in bosonic case, if $`P^{AB}`$ is invertible ($`P_{AB}P^{BC}=\delta _A^C`$), equation (1.114) is equivalent to the vanishing of the contorsion tensor $`H_{ABC}`$ defined by $$H_{ABC}=T_{AB}^DP_{DC}(1)^{ac}+T_{CA}^DP_{DB}(1)^{cb}+T_{BC}^DP_{DA}(1)^{ba}.$$ (1.116) This is the superspace generalization of (1.72), in terms of manifestly covariant quantities. The only nonvanishing components of $`H`$ are $`H_{\alpha \dot{\alpha }\beta }=i\left[P_{\alpha \dot{\alpha }\beta }+P_{\beta \dot{\alpha }\alpha }\right]`$ $`H_{\alpha \dot{\alpha }\dot{\beta }}=i\left[P_{\alpha \dot{\alpha }\dot{\beta }}+P_{\alpha \dot{\beta }\dot{\alpha }}\right]`$ $`H_{\alpha \dot{\alpha }\underset{¯}{b}}=iP_{\alpha \dot{\alpha }\underset{¯}{b}}`$ (1.117) We notice that its bosonic components $`H_{\underset{¯}{a}\underset{¯}{b}\underset{¯}{c}}`$ vanish due to the $`x`$–independence of the noncommutation functions in (1.98). The nonvanishing of $`H`$ comes entirely from the $`\theta `$–dependence of the functions in (1.100). As will be explained in section 1.3, it has been shown that bosonic coordinate-dependent deformations naturally emerge in string theory in the presence of curved backgrounds. Nonassociative deformations also emerge and the parameter governing nonassociativity is a bosonic three form $`H`$. Identifying the superspace analogue of this may help characterizing the superstring background where deformed superspaces may appear. When $`P`$ is invertible, equation (1.115) might seem to imply that the curvature $`R`$ is zero. This is not true in general because of the presence of the sign in front, which is dependent on the grading of the summed index $`M`$. Moreover, being puzzled by the unusual pattern of equations found, we have been trying to prove that (1.115) is algebraically implied by (1.114), but this doesn’t seem to work. We now search for the most general solutions of the conditions (1.114, 1.115). Writing them in terms of the $`P^{AB}`$ components we obtain $`B^{\alpha \dot{\beta }}A^{\beta \gamma }+A^{\beta \alpha }B^{\gamma \dot{\beta }}=0`$ $`B^{\alpha \dot{\beta }}B^{\beta \dot{\gamma }}+A^{\beta \alpha }\overline{A}^{\dot{\beta }\dot{\gamma }}=0`$ $`\left(\overline{C}^{\beta \dot{\beta }\dot{\alpha }}A^{\alpha \gamma }+C^{\beta \dot{\beta }\alpha }B^{\gamma \dot{\alpha }}\overline{C}^{\alpha \dot{\alpha }\dot{\beta }}A^{\beta \gamma }C^{\alpha \dot{\alpha }\beta }B^{\gamma \dot{\beta }}\right)=0`$ $`\mathrm{m}\left(\overline{C}^{\alpha \dot{\alpha }\dot{\beta }}C^{\gamma \dot{\gamma }\beta }+\overline{C}^{\gamma \dot{\gamma }\dot{\alpha }}C^{\beta \dot{\beta }\alpha }+\overline{C}^{\beta \dot{\beta }\dot{\gamma }}C^{\alpha \dot{\alpha }\gamma }\right)=0.`$ (1.118) The first two conditions necessarily imply the vanishing of the constants $`A`$ and $`B`$. Inserting this result in the third constraint we immediately realize that it is automatically satisfied and the only nontrivial condition which survives is the last one. This equation has nontrivial solutions. For example, the matrix $$C^{\alpha \dot{\alpha }\beta }=\left(\begin{array}{cc}\psi ^\beta & \psi ^\beta \\ \psi ^\beta & \psi ^\beta \end{array}\right)$$ (1.119) for any spinor $`\psi ^\beta `$, is a solution. It would correspond to assume the same commutations rules among any bosonic coordinate and the spinorial variables. We conclude that the most general associative and non(anti)commutative algebra in Minkowski superspace has the form $`\{\theta ^\alpha ,\theta ^\beta \}`$ $`=\{\overline{\theta }^{\dot{\alpha }},\overline{\theta }^{\dot{\beta }}\}=\{\theta ^\alpha ,\overline{\theta }^{\dot{\beta }}\}=0`$ $`[x^{\alpha \dot{\alpha }},\theta ^\beta ]`$ $`=iC^{\alpha \dot{\alpha }\beta }`$ $`[x^{\alpha \dot{\alpha }},\overline{\theta }^{\dot{\beta }}]`$ $`=i\overline{C}^{\alpha \dot{\alpha }\dot{\beta }}`$ (1.120) $`[x^{\alpha \dot{\alpha }},x^{\beta \dot{\beta }}]`$ $`=iD^{\alpha \dot{\alpha }\beta \dot{\beta }}+{\displaystyle \frac{1}{2}}\left(\overline{C}^{\beta \dot{\beta }\dot{\alpha }}\theta ^\alpha \overline{C}^{\alpha \dot{\alpha }\dot{\beta }}\theta ^\beta +C^{\beta \dot{\beta }\alpha }\overline{\theta }^{\dot{\alpha }}C^{\alpha \dot{\alpha }\beta }\overline{\theta }^{\dot{\beta }}\right),`$ where $`C`$ is subject to the last constraint in (1.118). Setting $`C^{\alpha \dot{\alpha }\beta }=0`$ we recover the usual noncommutative superspace considered so far in literature . Under conditions (1.118) the graded brackets (1.101) satisfy the Jacobi identities (1.113), as can be easily proved by expanding the functions in power series. In this case we have a well–defined super Poisson structure on superspace. We note that a non(anti)commutative but associative geometry always mantains the standard algebra (1.105) for the covariant derivatives. In fact, in this case, from (1.108) it follows $`R_{AB}^{}{}_{}{}^{CD}=0`$. ##### Construction of a Kontsevich-like product on superspace We will now describe the first few steps towards the construction of a star product defined on the class of general superfields. By definition, this product must be associative, i.e. it has to satisfy the Jacobi identities (1.113) when the fundamental algebra is associative. In section 1.1.4 we have seen that in the nonsupersymmetric case the lack of associativity of the fundamental algebra is signaled by the presence of a nonvanishing 3–form $`H`$. A product has been constructed so that the terms violating the Jacobi identities are proportional to $`H`$. The product is then automatically associative when the fundamental algebra is. In the present case we have shown that the lack of associativity in superspace is related to a nonvanishing super 3–form. This suggests the possibility to construct a super star product by suitably generalizing Kontsevich construction to superspace. The supersymmetric Poisson structure we have constructed on superspace can be written in a manifestly covariant form, as in (1.101), in terms of a constant matrix and covariant derivatives that have nontrivial torsion, or as in (1.103), in terms of a coordinate-dependent matrix and torsion free derivatives. The second formulation allows for a straightforward generalization of Kontsevich construction we outlined in section 1.1.4. However, the same procedure can be performed in a manifestly covariant way, and I choose to give this second version, so that the geometric interpretation of the results will be clear. We begin by considering Moyal–deformed product defined in the usual way $$\mathrm{\Phi }\mathrm{\Psi }\mathrm{\Phi }\mathrm{exp}(\mathrm{}\stackrel{}{D}_AP^{AB}\stackrel{}{D}_B)\mathrm{\Psi },$$ (1.121) where $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are arbitrary superfields, and $`\mathrm{}`$ denotes a deformation parameter. In general, due to the lack of (anti)commutativity among covariant derivatives (see eq. (1.106)), it is easy to prove that the $``$–product is not associative even when the Poisson brackets are. However, inspired by Kontsevich procedure , we perturbatively define a modified product $``$ with the property to be associative up to second order in $`\mathrm{}`$ when the Jacobi identities are satisfied. Precisely, we find an explicit form for the product by imposing the Jacobi identities (1.113) to be violated at this order only by terms proportional to $`H`$. To this end we define $`\mathrm{\Phi }\mathrm{\Psi }`$ $``$ $`\mathrm{\Phi }\mathrm{\Psi }+\mathrm{}\mathrm{\Phi }\stackrel{}{D}_AP^{AB}\stackrel{}{D}_B\mathrm{\Psi }+{\displaystyle \frac{\mathrm{}^2}{2}}\mathrm{\Phi }(\stackrel{}{D}_AP^{AB}\stackrel{}{D}_B)(\stackrel{}{D}_CP^{CD}\stackrel{}{D}_D)\mathrm{\Psi }`$ (1.122) $`{\displaystyle \frac{\mathrm{}^2}{3}}\left(\stackrel{}{D}_A\mathrm{\Phi }^{ABC}\stackrel{}{D}_B\stackrel{}{D}_C\mathrm{\Psi }(1)^c\stackrel{}{D}_C\stackrel{}{D}_A\mathrm{\Phi }^{ABC}\stackrel{}{D}_B\mathrm{\Psi }\right)`$ $`+𝒪(\mathrm{}^3),`$ where $`^{ABC}`$ $`=`$ $`P^{AD}T_{DE}^{}{}_{}{}^{C}P^{EB}(1)^{ce}+{\displaystyle \frac{1}{2}}P^{BD}T_{DE}^{}{}_{}{}^{A}P^{EC}(1)^{ae+a+b+ab+bc}`$ (1.123) $`+{\displaystyle \frac{1}{2}}P^{CD}T_{DE}^{}{}_{}{}^{B}P^{EA}(1)^{be+a+c+ac+ab}.`$ Since it is straightforward to show that $`(\mathrm{\Phi }\mathrm{\Psi })\mathrm{\Omega }\mathrm{\Phi }(\mathrm{\Psi }\mathrm{\Omega })=`$ (1.124) $`{\displaystyle \frac{2}{3}}\mathrm{}^2(1)^{(c+b)(e+1)+eg+cf}\stackrel{}{D}_A\mathrm{\Phi }P^{AE}P^{BF}P^{CG}H_{GFE}\stackrel{}{D}_C\mathrm{\Omega }\stackrel{}{D}_B\mathrm{\Psi }`$ $`+𝒪(\mathrm{}^3)`$ up to second order in $`\mathrm{}`$ the product is associative if and only if $`H=0`$, i.e. the fundamental algebra is associative. We note that at this order only the contorsion enters the breaking of associativity, being the curvature tensor $`R`$ of order $`\mathrm{}`$. In we did not pursue the construction of the star product to all orders in $`\mathrm{}`$ but we believed that in principle there were no obstructions to the generalization of Kontsevich procedure to all orders. The extension of the superspace product we proposed to all orders in $`\mathrm{}`$ has been obtained afterwards in . We now discuss the closure of the class of chiral superfields under the deformed products we have introduced. For a generic choice of the supermatrix $`P^{AB}`$ the star product of two chiral superfields (satisfying $`\overline{D}_{\dot{\alpha }}\mathrm{\Phi }=0`$) is not a chiral superfield, both for associative and nonassociative products. However, in the particular case where the only nonvanishing components of the symplectic supermatrix $`P^{AB}`$ are $`P^{\alpha \dot{\beta }}`$ and $`P^{\underset{¯}{a}\underset{¯}{b}}`$, chiral superfields are closed both under the deformed product defined in (1.121) and under Kontsevich star product (1.122) (for the latter up to terms of order $`𝒪(\mathrm{}^3)`$). Clearly for $`P^{\alpha \dot{\beta }}0`$ the above star products are no more associative. Because of this, it could be problematic to generalize for instance the Wess–Zumino model to non(anti)commutative superspace, since $`(\mathrm{\Phi }\mathrm{\Phi })\mathrm{\Phi }\mathrm{\Phi }(\mathrm{\Phi }\mathrm{\Phi })`$. However, one may notice that for chiral superfields the above star products become commutative<sup>3</sup><sup>3</sup>3Generalized star products that are commutative but nonassociative have been considered in a different context in .. This commutativity implies that there is no ambiguity in putting the parenthesis in the cubic interaction term of a deformed Wess–Zumino model, since, when $`\mathrm{\Phi }`$ is a chiral superfield, $`(\mathrm{\Phi }\mathrm{\Phi })\mathrm{\Phi }=\mathrm{\Phi }(\mathrm{\Phi }\mathrm{\Phi })`$ holds. Therefore, the action for the deformed Wess–Zumino model $$S=d^4xd^2\theta d^2\overline{\theta }\mathrm{\Phi }\overline{\mathrm{\Phi }}+d^4x\left[d^2\theta \left(\frac{m}{2}\mathrm{\Phi }\mathrm{\Phi }+\frac{g}{3}\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\right)+\text{c. c.}\right].$$ (1.125) is well defined and can be studied. Notice that in this case the $``$–product in the kinetic term cannot be simply substituted with the standard product as it happens in superspace geometries where grassmannian coordinates anticommute . ##### Non(anti)commutative $`N=2`$ Euclidean superspace From the discussion of $`d=4`$ $`N=1`$ superspace it is clear that the superspace conjugation relations relating $`\theta ^\alpha `$ and $`\overline{\theta }^{\dot{\alpha }}`$ are strong constraints on the non(anti)commutative algebra. The main difference in the description of euclidean superspace with respect to Minkowski relies on the reality conditions satisfied by the spinorial variables. So it is reasonable to hope that in an euclidean signature structures with a nonvanishing anticommutator between two fermions could appear. As it is well known , in euclidean signature a reality condition on spinors is applicable only in the presence of extended supersymmetry. In we concentrated on the simplest case, the $`N=2`$ euclidean superspace, even if our analysis can be easily extended to more general cases. In a chiral description the two–component Weyl spinors satisfy a symplectic Majorana condition $$(\theta _i^\alpha )^{}=\theta _\alpha ^iC^{ij}\theta _j^\beta C_{\beta \alpha },(\overline{\theta }^{\dot{\alpha },i})^{}=\overline{\theta }_{\dot{\alpha },i}\overline{\theta }^{\dot{\beta },j}C_{\dot{\beta }\dot{\alpha }}C_{ji}$$ (1.126) where $`C^{12}=C_{12}=i`$. This implies that the most general non(anti)commutative algebra can be written as an obvious generalization of (1.94) with the functions on the rhs now being in suitable representations of the R–symmetry group. When imposing covariance under (super)translations we obtain that the most general non(anti)commutative geometry in euclidean superspace is $`\{\theta _i^\alpha ,\theta _j^\beta \}`$ $`=A_{1}^{}{}_{ij}{}^{\alpha \beta ,},\{\overline{\theta }^{\dot{\alpha },i},\overline{\theta }^{\dot{\beta },j}\}=A_2^{\dot{\alpha }\dot{\beta },ij},\{\theta _i^\alpha ,\overline{\theta }^{\dot{\alpha },j}\}=B_i^{\alpha \dot{\alpha },j}`$ $`[x^{\alpha \dot{\alpha }},\theta _i^\beta ]`$ $`=i𝒞_{1}^{}{}_{i}{}^{\underset{¯}{a}\beta ,}(\theta ,\overline{\theta })`$ $`[x^{\alpha \dot{\alpha }},\overline{\theta }^{\dot{\beta },i}]`$ $`=i𝒞_2^{\alpha \dot{\alpha }\dot{\beta },i}(\theta ,\overline{\theta })`$ $`[x^{\alpha \dot{\alpha }},x^{\beta \dot{\beta }}]`$ $`=i𝒟^{\alpha \dot{\alpha }\beta \dot{\beta }}(\theta ,\overline{\theta })`$ (1.127) where $`𝒞_{1}^{}{}_{i}{}^{\alpha \dot{\alpha }\beta ,}(\theta ,\overline{\theta })C_{1}^{}{}_{i}{}^{\alpha \dot{\alpha }\beta ,}+{\displaystyle \frac{i}{2}}\theta _j^\alpha B_i^{\beta \dot{\alpha },j}+{\displaystyle \frac{i}{2}}\overline{\theta }^{\dot{\alpha },j}A_{1}^{}{}_{ji}{}^{\alpha \beta ,}`$ $`𝒞_2^{\alpha \dot{\alpha }\dot{\beta },i}(\theta ,\overline{\theta })C_2^{\alpha \dot{\alpha }\dot{\beta },i}+{\displaystyle \frac{i}{2}}\theta _j^\alpha A_2^{\dot{\alpha }\dot{\beta },ji}+{\displaystyle \frac{i}{2}}\overline{\theta }^{\dot{\alpha },j}B_j^{\alpha \dot{\beta },i}`$ $`𝒟^{\alpha \dot{\alpha }\beta \dot{\beta }}(\theta ,\overline{\theta })D^{\alpha \dot{\alpha }\beta \dot{\beta }}`$ $`+{\displaystyle \frac{1}{2}}\left(\theta _i^\alpha C_2^{\beta \dot{\beta }\dot{\alpha },i}\theta _i^\beta C_2^{\alpha \dot{\alpha }\dot{\beta },i}+\overline{\theta }^{\dot{\alpha },i}C_{1}^{}{}_{i}{}^{\beta \dot{\beta }\alpha ,}\overline{\theta }^{\dot{\beta },i}C_{1}^{}{}_{i}{}^{\alpha \dot{\alpha }\beta ,}\right)`$ $`+{\displaystyle \frac{i}{4}}\left(\theta _i^\alpha A_2^{\dot{\alpha }\dot{\beta },ij}\theta _j^\beta +\theta _i^\alpha B_j^{\beta \dot{\alpha },i}\overline{\theta }^{\dot{\beta },j}+\overline{\theta }^{\dot{\alpha },i}B_i^{\alpha \dot{\beta },j}\theta _j^\beta +\overline{\theta }^{\dot{\alpha },i}A_{1}^{}{}_{ij}{}^{\alpha \beta ,}\overline{\theta }^{\dot{\beta },j}\right)`$ (1.128) with $`A_1`$, $`A_2`$, $`B`$, $`C_1`$, $`C_2`$ and $`D`$ constant. Following the same steps as in the Minkowski case, we can look for the most general associative algebra. The results we obtain for associative non(anti)commuting geometries in euclidean superspace are $`\{\theta _i^\alpha ,\theta _j^\beta \}`$ $`=A_{1}^{}{}_{ij}{}^{\alpha \beta ,},\{\overline{\theta }^{\dot{\alpha },i},\overline{\theta }^{\dot{\beta },j}\}=0,\{\theta _i^\alpha ,\overline{\theta }^{\dot{\beta },j}\}=0`$ $`[x^{\alpha \dot{\alpha }},\theta _i^\beta ]`$ $`=iC_{1}^{}{}_{i}{}^{\alpha \dot{\alpha }\beta ,}{\displaystyle \frac{1}{2}}\overline{\theta }^{\dot{\alpha },j}A_{1}^{}{}_{ji}{}^{\alpha \beta ,}`$ $`[x^{\alpha \dot{\alpha }},\overline{\theta }^{\dot{\beta },i}]`$ $`=0`$ (1.129) $`[x^{\alpha \dot{\alpha }},x^{\beta \dot{\beta }}]`$ $`=iD^{\alpha \dot{\alpha }\beta \dot{\beta }}+{\displaystyle \frac{i}{2}}\left(\overline{\theta }^{\dot{\alpha },i}C_{1}^{}{}_{i}{}^{\beta \dot{\beta }\alpha ,}\overline{\theta }^{\dot{\beta },i}C_{1}^{}{}_{i}{}^{\alpha \dot{\alpha }\beta ,}\right){\displaystyle \frac{1}{4}}\overline{\theta }^{\dot{\alpha },i}A_{1}^{}{}_{ij}{}^{\alpha \beta ,}\overline{\theta }^{\dot{\beta },j}`$ or $`\{\theta _i^\alpha ,\theta _j^\beta \}`$ $`=0,\{\overline{\theta }^{\dot{\alpha },i},\overline{\theta }^{\dot{\beta },j}\}=A_2^{\dot{\alpha }\dot{\beta },ij},\{\theta _i^\alpha ,\overline{\theta }^{\dot{\beta },j}\}=0`$ $`[x^{\alpha \dot{\alpha }},\theta _i^\beta ]`$ $`=0`$ $`[x^{\alpha \dot{\alpha }},\overline{\theta }^{\dot{\beta },i}]`$ $`=iC_2^{\alpha \dot{\alpha }\dot{\beta },i}{\displaystyle \frac{1}{2}}\theta _j^\alpha A_2^{\dot{\alpha }\dot{\beta },ji}`$ (1.130) $`[x^{\alpha \dot{\alpha }},x^{\beta \dot{\beta }}]`$ $`=iD^{\alpha \dot{\alpha }\beta \dot{\beta }}+{\displaystyle \frac{i}{2}}\left(\theta _i^\alpha C_2^{\beta \dot{\beta }\dot{\alpha },i}\theta _i^\beta C_2^{\alpha \dot{\alpha }\dot{\beta },i}\right){\displaystyle \frac{1}{4}}\theta _i^\alpha A_2^{\dot{\alpha }\dot{\beta },ij}\theta _j^\beta `$ We notice that in this case associativity imposes less restrictive constraints because of the absence of conjugation relations between $`A_1`$ and $`A_2`$. As a consequence, nontrivial anticommutation relations among $`\theta `$’s (or $`\overline{\theta }`$’s) are allowed. Moreover the R–symmetry group of the $`N=2`$ euclidean superalgebra is broken only by the constant terms $`C_1`$ and $`C_2`$. Setting these terms equal to zero leads to nontrivial (anti)commutation relations preserving $`R`$–symmetry. Again, explicit expressions for the corresponding graded brackets can be obtained as an obvious generalization of (1.101, 1.102). In this case they define a super Poisson structure on euclidean superspace. A simple example of a super Poisson structure is $$\{\mathrm{\Phi },\mathrm{\Psi }\}_P=\mathrm{\Phi }\stackrel{}{D}_\alpha ^iA_{1}^{}{}_{ij}{}^{\alpha \beta ,}\stackrel{}{D}_\beta ^j\mathrm{\Psi }$$ (1.131) We notice that this extension is allowed only in euclidean superspace, where it is consistent with the reality conditions on the spinorial variables. #### 1.1.7 $`N=\frac{1}{2}`$ supersymmetry ##### Seiberg $`N=1/2`$ superspace In this section I will review nonanticommutative $`N=\frac{1}{2}`$ superspace introduced by Seiberg in and I will explain its relation with my results in . As I have discussed in detail in the last section, in I obtained the most general nonanticommutative deformation of $`d=4`$ $`N=1`$ superspace compatible with supersymmetry and associativity. This involves nontrivial $`[x,\theta ]`$, $`[x,\overline{\theta }]`$ and $`[x,x]`$, while $`\{\theta ,\theta \}`$, $`\{\overline{\theta },\overline{\theta }\}`$ and $`\{\theta ,\overline{\theta }\}`$ cannot be turned on. The situation improves in euclidean signature, where it is possible to turn on the fermionic anticommutators because of the different spinor conjugation relations. Rigorously, a superspace with euclidean signature can be defined only with extended supersymmetry, because of the impossibility to assign consistent reality conditions to $`\theta ^\alpha `$ and $`\overline{\theta }^{\dot{\alpha }}`$ in $`N=1`$ . This is why I have studied the $`N=2`$ euclidean superspace in . However, one can formally define an $`N=1`$ euclidean superspace by temporarily doubling the fermionic degrees of freedom, as it is done by Seiberg in . To show how this works, I will first redefine the $`N=2`$ euclidean spinor variables (1.126) as follows $`\theta ^\alpha \theta _1^\alpha \theta _2^\alpha ;\overline{\theta }^\alpha \theta _1^\alpha +\theta _2^\alpha `$ (1.132) $`\theta ^{\dot{\alpha }}\overline{\theta }^{1\dot{\alpha }}\overline{\theta }^{2\dot{\alpha }};\overline{\theta }^{\dot{\alpha }}\overline{\theta }^{1\dot{\alpha }}\overline{\theta }^{2\dot{\alpha }}`$ (1.133) These satisfy the reality conditions $`(\theta ^\alpha )^{}=i\overline{\theta }_\alpha ;(\overline{\theta }^\alpha )^{}=i\theta _\alpha `$ (1.134) $`(\theta _\alpha )^{}=i\overline{\theta }^\alpha ;(\overline{\theta }_\alpha )^{}=i\theta ^\alpha `$ (1.135) and analogous for the dotted variables. Dotted and undotted variables are unrelated. I will refer to $`\theta ^\alpha `$ and $`\theta ^{\dot{\alpha }}`$ as the left-moving sector of the theory and to $`\overline{\theta }^\alpha `$ and $`\overline{\theta }^{\dot{\alpha }}`$ as the right-moving sector. This terminology will be clarified in section 1.2.4, where the string theory origin of deformed superspaces will be discussed. We will see in section 1.2.4 that open string boundary conditions relate left- and right-moving fermionic variables on the D-brane ($`\theta =\overline{\theta }`$ on the boundary). As a result, the effective field theory on the brane is described by an $`N=1`$ euclidean superspace with fermionic variables $`\theta ^\alpha `$ and $`\theta ^{\dot{\alpha }}`$. For this reason from now on I will only consider spinor coordinates in the left-moving sector. In we have used in both Minkowski and euclidean signature a nonchiral representation for superspace covariant derivatives. As a result consistency with supersymmetry transformations implied that if we turned on a nonvanishing anticommutators between two fermionic variables, also nonzero boson-boson and boson-fermion commutators appeared. This made the algebra coordinate dependent and forced us to build a complicated Kontevich-like product. A way to obtain a simpler, constant algebra for the supercoordinates is the following . In four dimensions it is possible to use a chiral representation of supersymmetry, where, for each supersymmetry, one of the two covariant derivatives coincides with the ordinary one while the other gets dressed. If we make this chiral choice in $`N=2`$ euclidean superspace, by defining $`Q_\alpha =i(_\alpha i\theta ^{\dot{\alpha }}_{\alpha \dot{\alpha }});Q_{\dot{\alpha }}=i_{\dot{\alpha }}`$ (1.136) $`D_\alpha =_\alpha ;D_{\dot{\alpha }}=_{\dot{\alpha }}+i\theta ^\alpha _{\alpha \dot{\alpha }}`$ (1.137) corresponding to the supersymmetry transformations $$\delta x^{\alpha \dot{\alpha }}=iϵ^\alpha \theta ^{\dot{\alpha }};\delta \theta ^\alpha =ϵ^\alpha ;\delta \theta ^{\dot{\alpha }}=ϵ^{\dot{\alpha }}$$ (1.138) we immediately realize that the nonanticommutative algebra $$\{\theta ^\alpha ,\theta ^\beta \}=2P^{\alpha \beta }\mathrm{the}\mathrm{rest}=0$$ (1.139) with $`P^{\alpha \beta }`$ a constant symmetric matrix, is consistent with (1.138) and trivially associative . One can easily check that, comparing to (1.106, 1.107, 1.108), the supersymmetry algebra is not modified by the deformation but the algebra of covariant derivatives becomes $`\{D_\alpha ,D_\beta \}=0;\{D_\alpha ,D_{\dot{\alpha }}\}=i_{\alpha \dot{\alpha }}`$ (1.140) $`\{D_{\dot{\alpha }},D_{\dot{\beta }}\}=2P^{\alpha \beta }_{\alpha \dot{\alpha }}_{\beta \dot{\beta }}`$ (1.141) This is a problem when one tries to construct chiral superfields. So, even if this algebra is very simple, it is not suitable for constructing deformations of ordinary theories in superspace. Seiberg in chose the opposite chiral representation with respect to (1.137) $`Q_\alpha =i_\alpha ;Q_{\dot{\alpha }}=i(_{\dot{\alpha }}i\theta ^\alpha _{\alpha \dot{\alpha }})`$ (1.142) $`D_\alpha =_\alpha +i\theta ^{\dot{\alpha }}_{\alpha \dot{\alpha }};D_{\dot{\alpha }}=_{\dot{\alpha }}`$ (1.143) corresponding to the supersymmetry transformations $$\delta x^{\alpha \dot{\alpha }}=iϵ^{\dot{\alpha }}\theta ^\alpha ;\delta \theta ^\alpha =ϵ^\alpha ;\delta \theta ^{\dot{\alpha }}=ϵ^{\dot{\alpha }}$$ (1.144) When we turn on a nontrivial anticommutator between the $`\theta `$’s, we obtain $`\{\theta ^\alpha ,\theta ^\beta \}=2P^{\alpha \beta }\{\theta ^\alpha ,\theta ^{\dot{\beta }}\}=\{\theta ^{\dot{\alpha }},\theta ^{\dot{\beta }}\}=0`$ (1.145) $`[x^{\alpha \dot{\alpha }},\theta ^\beta ]=2iP^{\alpha \beta }\theta ^{\dot{\alpha }}`$ (1.146) $`[x^{\alpha \dot{\alpha }},x^{\beta \dot{\beta }}]=2\theta ^{\dot{\alpha }}P^{\alpha \beta }\theta ^{\dot{\beta }}`$ (1.147) This is consistent with supersymmetry and associativity (it is just the “$`N=1`$” version of the algebra (1.1.6) I gave in and discussed in the previous section). This is coordinate-dependent and would require the construction of a Kontsevich-like product for the superfield algebra. However, Seiberg observed that it is possible to make the change of variables $$y^{\alpha \dot{\alpha }}=x^{\alpha \dot{\alpha }}i\theta ^\alpha \theta ^{\dot{\alpha }}$$ (1.148) In terms of $`(y^{\alpha \dot{\alpha }},\theta ^\alpha ,\theta ^{\dot{\alpha }})`$ the superspace algebra takes the form (1.139) and this makes it possible to define a Moyal-like associative product acting on superfields $`\mathrm{\Phi }=\mathrm{\Phi }(y^{\alpha \dot{\alpha }},\theta ^\alpha ,\theta ^{\dot{\alpha }})`$ as follows $`\mathrm{\Phi }\mathrm{\Psi }`$ $`=`$ $`\mathrm{\Phi }\mathrm{exp}(\stackrel{}{}_\alpha P^{\alpha \beta }\stackrel{}{}_\beta )\mathrm{\Psi }`$ (1.149) $`=`$ $`\mathrm{\Phi }\mathrm{\Psi }\mathrm{\Phi }\stackrel{}{}_\alpha P^{\alpha \beta }\stackrel{}{}_\beta \mathrm{\Psi }{\displaystyle \frac{1}{2}}P^2^2\mathrm{\Phi }^2\mathrm{\Psi }`$ (1.150) where $`P^2P^{\alpha \beta }P_{\alpha \beta }`$. In contradistinction to the bosonic case, this Moyal-like product has a finite derivative expansion, because of the grassmannian nature of the variables involved in the deformation. With the antichiral representation (1.143) the covariant derivative algebra is not deformed, while the supersymmetry algebra is deformed by curvature terms analogous to the ones we found in (1.106, 1.107, 1.108) $`\{Q_\alpha ,Q_\beta \}=0;\{Q_\alpha ,Q_{\dot{\alpha }}\}=i_{\alpha \dot{\alpha }}`$ (1.151) $`\{Q_{\dot{\alpha }},Q_{\dot{\beta }}\}=2P^{\alpha \beta }_{\alpha \dot{\alpha }}_{\beta \dot{\beta }}`$ (1.152) Since only the dotted sector is modified, in its is argued that only $`N=\frac{1}{2}`$ is preserved. Exactly as in the $`N=2`$ case discussed in the previous section, these susy-breaking terms do not affect the supercoordinate algebra, that is consistent with supersymmetry. In chiral and vector superfields in $`N=\frac{1}{2}`$ have been extensively studied. In particular, since the covariant derivatives are not modified by the deformation, it is possible to define (anti)chiral superfields whose class is closed under $``$ and to write down the action for a deformed Wess-Zumino model $`S`$ $`=`$ $`{\displaystyle d^8z\overline{\mathrm{\Phi }}\mathrm{\Phi }}{\displaystyle \frac{m}{2}}{\displaystyle d^6z\mathrm{\Phi }^2}{\displaystyle \frac{\overline{m}}{2}}{\displaystyle d^6\overline{z}\overline{\mathrm{\Phi }}^2}`$ (1.153) $``$ $`{\displaystyle \frac{g}{3}}{\displaystyle d^6z\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }}{\displaystyle \frac{\overline{g}}{3}}{\displaystyle d^6\overline{z}\overline{\mathrm{\Phi }}\overline{\mathrm{\Phi }}\overline{\mathrm{\Phi }}}`$ (1.154) $`=`$ $`S(P=0)+{\displaystyle \frac{g}{6}}{\displaystyle d^4xP^2F^3}`$ (1.155) where $`F`$ is the auxiliary field in the chiral multiplet and total superspace derivatives have been neglected to obtain the last equality. Similarly, one can see that the $`N=\frac{1}{2}`$ deformation of super Yang-Mills is characterized by explicitly susy-breaking P-dependent component terms . Moreover in the antichiral ring defined by the operator relation $`[Q_\alpha ,\overline{𝒪}]=0`$ has been studied and it has been shown that all its properties are preserved by the deformation, while the chiral ring cannot be defined since $`Q_{\dot{\alpha }}`$ is not a symmetry of the theory anymore. ##### Results in non(anti)commutative field theories After Seiberg’s paper appeared, deformed superspaces have attracted much attention and a lot of efforts have been done to elucidate the properties of nonanticommutative field theories in superspace. This interest is mostly due to the fact that nonanticommutative superspaces have been shown to naturally emerge in superstring theory in the presence of R-R backgrounds . I will review the string theory side of the story in section 1.2.4. Here I’m going to browse the huge bibliography and discuss the main results obtained, without giving any detail. I will first discuss the progress in understanding $`N=\frac{1}{2}`$ theories. In the $`N=\frac{1}{2}`$ deformation of WZ and super Yang-Mills theories have been proposed. The deformation of WZ model is easy to describe, since in component formulation it corresponds to adding to the ordinary action a cubic term in the auxiliary field $`F`$. In some features of the deformed WZ model, such as non validity of standard nonrenormalization theorems, stability of the vacuum energy and existence of the antichiral ring have been discussed through some examples, in both component and superspace formulations. In particular, since supersymmetry plays an important role to guarantee renormalization through partial cancellation of UV divergences associated to bosonic and fermionic degrees of freedom, it is compelling to study renormalizability properties of the $`N=\frac{1}{2}`$ theory where part of the supersymmetry is explicitly broken. A systematical analysis of perturbative renormalizability of $`N=\frac{1}{2}`$ WZ model has been performed in by explicit calculations up to two loops. It has been shown that, even if new divergences appear, the model can be rendered renormalizable by adding ab initio $`F`$ and $`F^2`$ terms to the ordinary lagrangian. It is somehow expected that these terms may accompany the $`F^3`$ deformation, since they are allowed by the symmetry of the theory. These two-loop results have been extended to all orders in perturbation theory in both component and superspace formulations . The proof of renormalizability has been given on the base of dimensional arguments and global symmetries. The study of deformations of gauge theories is more interesting than the scalar case. $`N=\frac{1}{2}`$ $`U(n)`$ gauge theories have been proven to be renormalizable in in WZ gauge. This result have been checked up to one loop in in component formulation, again in WZ gauge. The study of these gauge theories in a manifestly gauge independent superspace setup has been accomplished by generalizing the background field method to nonanticommutative case . In $`N=\frac{1}{2}`$ super Yang-Mills instantons have been studied. In the problem of constructing a Seiberg-Witten map analogous to the one discussed in section 1.2.1 for superfields in deformed superspaces has been considered. The case with extended supersymmetry has also been considered and it has been shown that, while in general $`N=(1,1)`$ supersymmetry is broken to $`N=(\frac{1}{2},0)`$, there are particular cases where $`N=(\frac{1}{2},1)`$ survives. Moreover, super Yang-Mills theory on extended deformed superspaces has also been studied in . $`d=2`$ $`N=2`$ classical aspects of sigma models characterized by a general Kähler potential and arbitrary superpotential deformed by a nonanticommutative product have been studied in . Finally, the connection between nonanticommutative geometry and supermatrix models have been studied in . ### 1.2 Non(anti)commutative field theory from the (super)string #### 1.2.1 Noncommutative Yang-Mills theory from the open string In this section I mostly refer to and show that a constant Neveu-Schwarz Neveu-Schwarz (NS-NS) $`B`$-field background modifies string dynamics nontrivially when D-branes are present. The open string with extrema constrained to lay on Dp-branes sees a deformed target-space $`G_{\mu \nu }`$ metric and a noncommutative target space coordinate algebra, characterized by a matrix $`\theta ^{\mu \nu }`$. Taking a zero slope limit $`\alpha ^{}0`$ in such a way to keep the above open-string parameters finite, one can obtain two different effective theory descriptions. Depending on the choice of the regularization prescription, one obtains a field theory in ordinary space where the background field appears explicitly, or a noncommutative field theory where the background field only appears implicitly in the noncommutativity matrix $`\theta ^{\mu \nu }`$. We are not going to give a review of the basic string theory needed in this section. For this we suggest the textbooks . ##### The open string effective metric and noncommutativity parameter Let us consider the bosonic sector of open string theory, in a 10-dimensional flat spacetime background with metric $`g_{\mu \nu }`$, in the presence of a constant NS-NS field $`B_{\mu \nu }`$ and Dp-branes. Let us assume $`B_{0i}=0`$, where $`i`$ is a generic spacelike direction and $`0`$ is timelike. This means that we are going to consider a magnetic $`B`$ field. At the end of this section I will briefly discuss the electric case, to show that a zero-slope limit giving a noncommutative effective field theory on the brane is not admitted in this case . It is well-known that a constant background $`B`$ field can be gauged away in the bulk, but not on the boundary, on the Dp-brane, where it acts as a constant magnetic (or electric) field . If $`\mathrm{rk}(B)=r`$, we can assume $`rp+1`$. We will choose spacetime coordinates in such a way that $`B_{ij}0`$ for $`i,j=1,\mathrm{},r`$ only and $`g_{ij}=0`$ for $`i=1,\mathrm{},r,j1,\mathrm{},r`$. The worldsheet action is $`S={\displaystyle \frac{1}{4\pi \alpha ^{}}}{\displaystyle _\mathrm{\Sigma }}\left(g_{\mu \nu }_ax^\mu ^ax^\nu 2\pi i\alpha ^{}B_{ij}ϵ^{ab}_ax^i_bx^j\right)`$ (1.156) $`={\displaystyle \frac{1}{4\pi \alpha ^{}}}{\displaystyle _\mathrm{\Sigma }}g_{\mu \nu }_ax^\mu ^ax^\nu {\displaystyle \frac{i}{2}}{\displaystyle _\mathrm{\Sigma }}B_{ij}x^i_tx^j`$ (1.157) where the second equality shows that antisymmetry of $`ϵ`$ implies that the $`B`$-term is a boundary term. $`\mathrm{\Sigma }`$ is the (euclidean) string worldsheet, $`\mathrm{\Sigma }`$ is its boundary, $`_t`$ is the derivative tangent to the boundary $`\mathrm{\Sigma }`$. Because of the presence of Dp-branes, the boundary term cannot be eliminated. It modifies the boundary conditions for the open string in the directions $`i`$ along the brane $$g_{ij}_nx^j+2\pi i\alpha ^{}B_{ij}_tx^j\|_\mathrm{\Sigma }=0$$ (1.158) where $`_n`$ is the derivative in the normal direction with respect to the boundary $`\mathrm{\Sigma }`$. In this equation both i and j are along the brane. We observe that for $`B=0`$ we obtain Neumann boundary conditions, while for maximal $`B`$ rank on the brane and $`B\mathrm{}`$ we get Dirichlet ones. In the latter case the string extrema are constrained to a single point on the Dp-brane, since the coordinates along the brane that describe them do not move. We can thus think that string extrema are attached to a 0-brane on the p-brane. From now on we will consider the classical approximation to string theory. $`\mathrm{\Sigma }`$ is a disk, that can be mapped to the upper half plane, described by complex coordinates $`z`$, $`\overline{z}`$, because of conformal invariance. The boundary conditions (1.158) can then be rewritten as $$g_{ij}(\overline{})x^j+2\pi \alpha ^{}B_{ij}(+\overline{})x^j\|_{z=\overline{z}}=0$$ (1.159) where $`=\frac{}{z}`$, $`\overline{}=\frac{}{\overline{z}}`$, $`Im(z)0`$. The propagator $`x^i(z)x^j(z^{})`$ with the boundary conditions (1.159) is given by $`x^i(z)x^j(z^{})=\alpha ^{}[g^{ij}\mathrm{log}\|zz^{}\|g^{ij}\mathrm{log}\|z\overline{z}^{}\|`$ (1.160) $`+G^{ij}\mathrm{log}\|z\overline{z}^{}\|^2+{\displaystyle \frac{1}{2\pi \alpha ^{}}}\theta ^{ij}\mathrm{log}{\displaystyle \frac{z\overline{z}^{}}{\overline{z}z^{}}}+D^{ij}]`$ (1.161) where $`G^{ij}=\left({\displaystyle \frac{1}{g+2\pi \alpha ^{}B}}\right)_S^{ij}=\left({\displaystyle \frac{1}{g+2\pi \alpha ^{}B}}g{\displaystyle \frac{1}{g2\pi \alpha ^{}B}}\right)^{ij}`$ (1.162) $`G_{ij}=g_{ij}\left(2\pi \alpha ^{}\right)^2\left(Bg^1B\right)_{ij}`$ (1.163) $`\theta ^{ij}=2\pi \alpha ^{}\left({\displaystyle \frac{1}{g+2\pi \alpha ^{}B}}\right)_A^{ij}=\left(2\pi \alpha ^{}\right)^2\left({\displaystyle \frac{1}{g+2\pi \alpha ^{}B}}B{\displaystyle \frac{1}{g2\pi \alpha ^{}B}}\right)^{ij}`$ (1.164) and $`()_S`$, $`()_A`$ denote the symmetric and antisymmetric part of the matrix in brackets, respectively. The constant quantities $`D^{ij}`$ depend on $`B`$, but not on $`z`$ and $`z^{}`$. They can be fixed to a certain value by making use of the fact that $`B`$ is arbitrary. We are interested in taking the limit $`z\tau R`$, $`z^{}\tau ^{}R`$ in (1.161). This is because in the open string case vertex operators are to be inserted on the boundary of the worldsheet. Taking the limit one obtains $$x^i(\tau )x^j(\tau ^{})=\alpha ^{}G^{ij}\mathrm{log}(\tau \tau ^{})^2+\frac{i}{2}\theta ^{ij}ϵ(\tau \tau ^{})$$ (1.165) The function $`ϵ(\tau )`$ is $`1`$ ($`1`$) for positive (negative) $`\tau `$. The discontinuity in the propagator can be expressed in terms of the function $`ϵ`$ when convenient values for the constants $`D^{ij}`$ are chosen. $`G^{ij}`$ is interpreted as the effective metric seen by open strings. This can be understood by comparing with the closed string case, where the propagator between two internal worldsheet points has a short distance behavior given by $$x^i(z)x^j(z^{})=\alpha ^{}g^{ij}\mathrm{log}\|zz^{}\|$$ (1.166) It is clear from (1.165) that, in the commutative limit $`\theta 0`$, the metric $`G^{ij}`$ is for the open string what $`g^{ij}`$ is for the closed string. By considering the second term in (1.165), we can see that the coefficient $`\theta ^{ij}`$ can be interpreted as noncommutativity parameter for the coordinates along the Dp-brane. In conformal field theory there exists a correspondence between time ordering and operator ordering, that in this case gives $$[x^i(\tau ),x^j(\tau )]=T\left(x^i(\tau )x^j(\tau ^{})x^i(\tau )x^j(\tau ^+)\right)=i\theta ^{ij}$$ (1.167) The first equality says that the path integral of the time-ordered combination in the right hand side corresponds to a matrix element of the equal-time commutator between $`x^i`$ and $`x^j`$. So we deduce that the coordinates $`x^i`$ are noncommuting with parameter $`\theta ^{ij}`$. ##### Correlation functions, effective action and Moyal product Let us now consider the product of two tachyon vertex operators $`e^{ipx}(\tau )`$, $`e^{iqx}(\tau ^{})`$, with $`\tau >\tau ^{}`$. By contracting with the two point function (1.165) we obtain the short distance behavior $`e^{ipx}(\tau )e^{iqx}(\tau ^{})(\tau \tau ^{})^{2\alpha ^{}G^{ij}p_iq_j}e^{\frac{1}{2}i\theta ^{ij}p_iq_j}e^{i(p+q)x}(\tau ^{})+\mathrm{}`$ (1.168) We observe that in the limit $`\alpha ^{}0`$ the OPE formula reduces to $``$ product $`e^{ipx}(\tau )e^{iqx}(\tau ^{})e^{ipx}e^{iqx}(\tau ^{})`$ (1.169) Therefore we expect that in the limit $`\alpha ^{}0`$ the theory should be easily described in terms of Moyal product. However, many interesting aspects emerge even without taking that limit. We have shown that open strings in the presence of a constant $`B`$ field on a Dp-brane can be described not only in terms of the two parameters $`g^{ij}`$ and $`B^{ij}`$, but also in terms of two functions $`G^{ij}(g,B;\alpha ^{})`$ and $`\theta ^{ij}(g,B;\alpha ^{})`$, representing the effective metric and noncommutativity parameter seen by the open string ending on the brane. Now we are going to show that, if we choose to work with the second couple of parameters, the dependence of the effective theory on $`\theta `$ is very simple. The theory with a nonvanishing $`\theta `$ can be obtained from the one with $`\theta =0`$ by simply substituting Moyal $``$ products characterized by the noncommutativity parameter $`\theta `$ to ordinary products. This is a general feature, no limit has to be taken. It is well-known that perturbative string theory requires the evaluation of the path integral of vertex operators the generally take the form $`P(x,^2x,\mathrm{})e^{ipx}`$, where $`P`$ is a polynomial in the derivatives of $`x`$. Let us now consider the expectation value of the product of $`k`$ vertex operators, with momenta $`p^1,\mathrm{},p^k`$ and $`x`$ along the Dp-brane. We are interested in determining the explicit dependence on the parameter $`\theta `$, while we would like to keep the dependence on $`G`$ implicit. The two-point function is the sum of two terms, one contains $`G`$ only, the other contains $`\theta `$ only. When we contract terms containing derivatives of $`x`$, the part of the propagator involving $`\theta `$ does not contribute, because $`ϵ(\tau \tau ^{})`$ is a constant function. So we can obtain the $`\theta `$ dependence by just taking into consideration the contraction between exponentials $`{\displaystyle \underset{n=1}{\overset{k}{}}}P_n(x(\tau _n),^2x(\tau _n),\mathrm{})e^{ip^nx(\tau _n)}_{G,\theta }`$ (1.170) $`=e^{\frac{i}{2}_{n>m}p_i^n\theta ^{ij}p_j^mϵ(\tau _n\tau _m)}{\displaystyle \underset{n=1}{\overset{k}{}}}P_n(x(\tau _n),^2x(\tau _n),\mathrm{})e^{ip^nx(\tau _n)}_{G,\theta =0}`$ (1.171) (1.172) So the $`\theta `$-dependence is simply given by the phase factor $$e^{\frac{i}{2}_{n>m}p_i^n\theta ^{ij}p_j^mϵ(\tau _n\tau _m)}$$ (1.173) Because of momentum conservation $`_np^n=0`$ and antisymmetry of $`\theta `$, (1.173) only depends on the cyclic order of the points $`\tau _1,\mathrm{},\tau _n`$ on the worldsheet boundary. By knowing the $`S`$-matrix of low energy massless particles, one can deduce order by order in $`\alpha ^{}`$ an effective action for the theory. This will be expressed in terms of a certain number of functions $`\mathrm{\Phi }_i`$ that in general may have values in the space of $`n\times n`$ matrices (think of the nonabelian gauge theory case). $`\mathrm{\Phi }_i`$ represents the wave function of the i-th field. By looking at (1.164) one notes that $`B=0\theta =0`$. So the general form of the effective action for $`B=0`$ will be as follows $$d^{p+1}x\sqrt{\mathrm{det}G}Tr\left(^{n_1}\mathrm{\Phi }_1^{n_2}\mathrm{\Phi }_2\mathrm{}^{n_k}\mathrm{\Phi }_k\right)$$ (1.174) where $`^{n_i}`$ stands for the product of $`n_i`$ partial derivatives with respect to certain unspecified coordinates. Now it is easy to move on to a theory with $`B0`$. If the effective action is written in momentum space, then it is sufficient to insert the phase factor (1.173). In configuration space this corresponds to replacing ordinary products with Moyal $``$ products (see (1.36)). To conclude, the effective action with $`B0`$ takes the form $$d^{p+1}x\sqrt{\mathrm{det}G}Tr\left(^{n_1}\mathrm{\Phi }_1^{n_2}\mathrm{\Phi }_2\mathrm{}^{n_k}\mathrm{\Phi }_k\right)$$ (1.175) So we have found an easy way to describe the theory with $`B0`$ when knowing the one with $`B=0`$. However we must stress that both theories have an equally complicated $`\alpha ^{}`$ expansion. ##### The description in the zero-slope limit The formalism of noncommutative geometry becomes much more powerful when the zero-slope limit $`\alpha ^{}0`$ is taken. This is somehow expected from (1.169). We would like to take the limit in such a way to keep the open string parameters $`G`$ and $`\theta `$ finite. This can be done by choosing $`\alpha ^{}ϵ^{\frac{1}{2}}0`$ (1.176) $`g_{ij}ϵ0\mathrm{per}\mathrm{i},\mathrm{j}=1,\mathrm{},\mathrm{r}`$ (1.177) when all the rest is kept fixed (also the two form $`B`$). Equation (1.164) becomes $`G^{ij}=\{\begin{array}{cc}\frac{1}{(2\pi \alpha ^{})^2}\left(\frac{1}{B}g\frac{1}{B}\right)^{ij}& \mathrm{per}\mathrm{i},\mathrm{j}=1,\mathrm{},\mathrm{r}\\ g^{ij}& \mathrm{otherwise}\end{array}`$ (1.178) $`G_{ij}=\{\begin{array}{cc}(2\pi \alpha ^{})^2(Bg^1B)_{ij}& \mathrm{per}\mathrm{i},\mathrm{j}=1,\mathrm{},\mathrm{r}\\ g_{ij}& \mathrm{otherwise}\end{array}`$ (1.179) $`\theta ^{ij}=\{\begin{array}{cc}\left(\frac{1}{B}\right)^{ij}& \mathrm{per}\mathrm{i},\mathrm{j}=1,\mathrm{},\mathrm{r}\\ 0& \mathrm{otherwise}\end{array}`$ (1.180) The propagator for two points on the boundary becomes $$x^i(\tau )x^j(0)=\frac{i}{2}\theta ^{ij}ϵ(\tau )$$ (1.181) For two generic functions $`f`$ and $`g`$ we then obtain (see (1.169)) $$:f(x(\tau ))::g(x(0)):=:e^{\frac{i}{2}ϵ(\tau )\theta ^{ij}\frac{}{x^i(\tau )}\frac{}{x^j(0)}}f(x(\tau ))g(x(0)):$$ (1.182) and thus $$\underset{\tau 0^+}{lim}:f(x(\tau ))::g(x(0)):=:f(x(0))g(x(0)):$$ (1.183) where $``$ is Moyal product (1.32). As a result, correlation functions of exponential operators on the disk boundary are given by $$\underset{n}{}e^{ip_i^nx^i(\tau _n)}=e^{\frac{i}{2}_{n>m}p_i^n\theta ^{ij}p_j^mϵ(\tau _n\tau _m)}\delta \left(p^n\right)$$ (1.184) In the general case with $`n`$ functions $`f_1,\mathrm{},f_n`$ $$\underset{n}{}f_n(x(\tau _n))=𝑑xf_1(x)\mathrm{}f_n(x)$$ (1.185) ##### Adding gauge fields Let us now add to the action (1.157) a term representing the coupling of the string worldsheet to a gauge field $`A_i(x)`$. For simplicity we will take $`\mathrm{rk}(A)=1`$. $$i𝑑\tau A_i(x)_\tau x^i$$ (1.186) Comparing with (1.157) we see that the constant field $`B`$ can be replaced by a gauge field $`A_i=\frac{i}{2}B_{ij}x^j`$, whose field strength is $`F=B`$. The bosonic string coupled to the $`B`$-field background is invariant under the gauge symmetry $$\delta B_{\mu \nu }=_{[\mu }\mathrm{\Lambda }_{\nu ]}$$ (1.187) modulo boundary terms. These terms can be compensated by the shift $$\delta A_\mu =\mathrm{\Lambda }_\mu $$ (1.188) Therefore, physics will be described by the gauge invariant combination $`\omega =F+B`$. The action (1.186) is invariant under the transformation $$\delta A_i=_i\lambda $$ (1.189) since the variation of the integrand is a total derivative. When we consider a quantum field theory we must pay attention to the regularization procedure. Physics of course must not depend on the particular choice of regularization! If we choose a Pauli-Villars regularization, we obtain an ordinary gauge theory, symmetric under the usual gauge transformation (1.189). We can make a different choice, though. We can use a point-splitting regularization, characterized by the fact that the product of two operators at the same point never appears. Actually, one first eliminates the region $`\|\tau \tau ^{}\|<\delta `$ and then takes the limit $`\delta 0`$. Now we are going to evaluate the variation of the path integral of the exponential of the action (1.186) under the gauge transformation (1.189), having first expanded the exponential in series with respect to $`A`$. The first term in the expansion gives $$𝑑\tau A_i(x)_\tau x^i𝑑\tau ^{}_\tau ^{}\lambda $$ (1.190) Even though the integrand of the second factor is a total derivative, one gets a boundary contribution for $`\tau \tau ^{}=\pm \delta `$. In the limit $`\delta 0`$ this contribution takes the form $`{\displaystyle }d\tau :A_i(x(\tau ))_\tau x^i(\tau )::(\lambda (x(\tau ^{}))\lambda (x(\tau ^+))):`$ (1.191) $`={\displaystyle 𝑑\tau }:\left(A_i(x)\lambda \lambda A_i(x)\right)_\tau x^i:`$ (1.192) To obtain this results one makes use of the fact that there are no contractions between $`_\tau x`$ and $`x`$ with the constant propagator (1.181). Moreover, the relation (1.183) between operators and $``$ product has been taken into account. To cancel out the term (1.192), we have to modify the transformation (1.189). So we discover that the point splitting regularized theory is not invariant under (1.189), but under the new transformation $$\delta \widehat{A}_i=_i\widehat{\lambda }+i\widehat{\lambda }\widehat{A}_ii\widehat{A}_i\widehat{\lambda }$$ (1.193) We recognize the gauge invariance (1.54) of noncommutative gauge theories with $`n=1`$. We have introduced the “hatted” notation $`\widehat{A}`$ for fields in a noncommutative algebra. It is possible to show that at a generic order $`m`$ in the $`A`$-expansion of the exponential of the action the correct gauge invariance is the noncommutative one, when point-splitting regularization is performed. Moreover in the calculation of the expectation value of three gauge vertex operators is performed. The result of this computation could also be obtained by considering the following effective action as a starting point $$S_{\mathrm{eff}}\sqrt{\mathrm{det}G}G^{ii^{}}G^{jj^{}}Tr\left(\widehat{F}_{ij}\widehat{F}_{i^{}j^{}}\right)$$ (1.194) where the field strength can be expressed in terms of $`\widehat{A}`$ as in (1.54). We have shown that the same string theory in the limit $`\alpha ^{}0`$ can be either described by an effective action corresponding to an ordinary gauge theory or by one associated to a noncommutative gauge theory, depending on the choice of the regularization procedure. Since physics cannot depend on the way this procedure is performed, two results obtained with different regularizations must be related by a redefinition of the coupling constants. In the worldsheet action the spacetime-valued fields play the role of coupling constants, so we expect that commutative and noncommutative effective descriptions should be related by a redefinition of these fields. A natural guess is that a local map $`\widehat{A}=\widehat{A}(A,A,^2A,\mathrm{};\theta )`$ among gauge fields and a corresponding map $`\widehat{\lambda }=\widehat{\lambda }(\lambda ,\lambda ,^2\lambda ,\mathrm{};\theta )`$ among the group parameters exists. In section 1.1.2 we observed that noncommutative $`U(1)`$ group is nonabelian. This tells us that such a correspondence cannot exist. Actually, the existence of such a map would imply an isomorphism between the ordinary gauge group and the corresponding noncommutative one. Since an abelian group cannot be isomorphic to a non abelian one, the first proposal for the map is ruled out. However, what is really needed for physics is that the gauge-transformed field $`\delta _\lambda A`$ corresponds to the gauge-transformed field $`\widehat{\delta }_{\widehat{\lambda }}\widehat{A}`$. Therefore it is sufficient that $$\widehat{A}(A)+\widehat{\delta }_{\widehat{\lambda }}\widehat{A}(A)=\widehat{A}(A+\delta _\lambda A)$$ (1.195) and we can look for a correspondence $`\widehat{A}=\widehat{A}(A)`$ (1.196) $`\widehat{\lambda }=\widehat{\lambda }(\lambda ,A)`$ (1.197) satisfying (1.195). The $`A`$-dependence of $`\widehat{\lambda }`$ solves the problem of the isomorphism between the two gauge groups. A relation like (1.197) does not imply any correspondence between the two group structures. A correspondence of the required form was found in (Seiberg-Witten map). It is given in terms of a set of differential equations describing how $`A`$ and $`\lambda `$ must vary with $`\theta `$ for the physics to remain unchanged. The results discussed in this section concerning an open string ending on a single Dp-brane can be easily generalized to the case of a stack of $`n`$ coincident D-branes. In this case one obtains a $`U(n)`$ noncommutative Yang-Mills theory as an effective field theory on the brane worldvolume. In section 1.2.2 I have commented on the fact that $`SU(n)`$, $`SO(n)`$ and $`Sp(n)`$ subgroups are not closed under Moyal product, so in principle one expects that restrictions to this subgroups should not emerge from the string. Actually, the case of $`SU(n)`$ is ruled out because the $`U(1)`$ degree of freedom in $`U(n)`$ ceases to decouple. $`SO(n)`$ and $`Sp(n)`$ restrictions can instead emerge, in a very nontrivial way, by orientifold projection, as shown in . ##### Open string in the presence of an electric NS-NS background and the breakdown of unitarity and causality in time-space noncommutative field theories In section 1.1.2 I have discussed unitarity and causality problems in noncommutative field theories with time-space noncommutativity, i.e. with $`\theta ^{0i}0`$. In principle we could expect these theories to arise in string theory in the presence of D-branes and a constant electric $`B_{0i}`$ background. However, the case of an electric background field is very different with respect to the magnetic one. It can be shown that if the background electric field $`E`$ exceeds the critical upper value $`E_c`$, string pairs are produced that destabilize the vacuum. So, if the electric field is along the $`x_1`$ direction and the metric is diagonal in the $`(x_0,x_1)`$ plane with components given by $`g`$, the bound is given by $$EE_c,\mathrm{where}E_c=\frac{g}{2\pi \alpha ^{}}$$ (1.198) In it has been shown that in this case the open string parameters $`(G,\theta )`$ are related to the electric field on the brane by the formula $$\alpha ^{}G^1=\frac{1}{2\pi }\frac{E}{E_c}\theta $$ (1.199) As before, to obtain the effective field theory on the brane we have to consider a zero-slope limit $`\alpha ^{}0`$. From the previous formula, it is clear that if we want to keep the open string metric $`G`$ finite in the limit, when $`\alpha ^{}0`$ then also $`\theta 0`$. Therefore, it is possible to obtain a field theory description involving massless open string modes only, but this will be an ordinary field theory and not a noncommutative one. On the other hand, we can keep the noncommutativity parameter $`\theta `$ finite, but then $`\alpha ^{}`$ must also be finite and we are considering a string theory and not a field theory. Indeed, in , it has been shown that a limit can be taken where $$\frac{E}{E_c}1\mathrm{and}g\frac{1}{1\left(\frac{E}{E_c}\right)^2}$$ (1.200) and all the other parameters are kept fixed, in particular the open string metric $`G`$. In this limit $$\theta =2\pi \alpha ^{}G^1$$ (1.201) is finite, so time-space noncommutativity is present. The theory obtained describes open strings in noncommutive spacetime (NCOS). Open strings decouple from closed strings, therefore also from gravity. In it has been shown that NCOS theory is the S-dual description of strongly coupled, spatially noncommutative $`N=4`$ Yang-Mills theory (other works concerning NCOS are listed in ). We conclude that noncommutative field theory with time-space noncommutativity does not emerge as a consistent truncation of string theory. Moreover, string theory in the presence of an electric background on the brane is unitary and acausal effects are not present . So $`\alpha ^{}`$ corrections to noncommutative “electric” field theory restore unitarity and causality. It is clear that “electric” noncommutative field theories are missing some degrees of freedom, related to the undecoupled massive string modes, that are necessary for unitarity and causality of the theory. In it has been shown that tachyonic particles are produced in scattering processes of noncommutative field theory with $`\theta ^{0i}`$ turned on. Form the string theory point of view, these particles may be viewed as a remnant of a continuous spectrum of these undecoupled closed-string modes. #### 1.2.2 Generalization to the superstring in RNS and GS formalisms In this section I will generalize to the superstring the results obtained in the previous section for the bosonic string. In the string with $`N=1`$ worldsheet supersymmetry (RNS) is considered. In instead the manifestly target-space supersymmetric superstring (GS) is discussed. In both cases the open superstring is coupled to a constant NS-NS background in the presence of D-branes. I’m not going to give an introduction to these two formalisms for the superstring. For this I suggest the textbooks ##### Open RNS string in the presence of a constant $`B`$-field and D-branes In the following action for the RNS string coupled to a constant magnetic NS-NS background field $`B^{ij}`$ is considered $`S={\displaystyle \frac{1}{4\pi \alpha ^{}}}{\displaystyle d^2z\left\{\overline{}x^\mu x_\mu +i\psi ^\mu \overline{}\psi _\mu +i\overline{\psi }^\mu \overline{\psi }_\mu 2\pi i\alpha ^{}B_{ij}ϵ^{ab}_ax^i_bx^j\right\}}`$ (1.202) (1.203) In the directions $`i`$, $`j`$ along the brane the following boundary conditions are imposed $`g_{ij}(\overline{})x^j+2\pi \alpha ^{}B_{ij}(+\overline{})x^j\|_{z=\overline{z}}=0`$ (1.204) $`g_{ij}(\psi ^j\overline{\psi }^j)+2\pi \alpha ^{}B_{ij}(\psi ^j+\overline{\psi }^j)\|_{z=\overline{z}}=0`$ (1.205) The first condition can be naturally obtained by requiring that there are no boundary terms in the Euler-Lagrange equations of motion. The second one, involving fermions, is obtained by requiring consistency under supersymmetry of the boundary conditions, but cannot be obtained from the action (1.203). Actually, one finds an inconsistency when requiring both the vanishing of boundary terms in the variation of the action and the compatibility of boundary conditions with supersymmetry. In (see also for a nice summary of the methods used) it has been shown that boundary terms can be added to the action (1.203) so that the supersymmetric boundary conditions (1.205) follow as boundary contributions to the field equations. Therefore, the theory described by the action $`S+S_b`$, with $$S_b=\frac{1}{2}d^2zB_{\mu \nu }\overline{\psi }^\mu \rho ^\alpha _\alpha \psi ^\nu $$ (1.206) and by the boundary conditions (1.205) is invariant under the rigid $`N=1`$ worldsheet supersymmetry transformations $`\delta x^i=i\eta (\psi ^i+\overline{\psi }^i)`$ (1.207) $`\delta \psi ^i=\eta x^i`$ (1.208) $`\delta \overline{\psi }^i=\eta \overline{}x^i`$ (1.209) where the parameter $`\eta `$ is a worldsheet spinor and a spacetime scalar. In the open string is coupled to a gauge field $`A`$ by adding to (1.203) the following boundary term $$L_A=i𝑑\tau \left(A_i(x)_\tau x^iiF_{ij}\mathrm{\Psi }^i\mathrm{\Psi }^j\right)$$ (1.210) where $`F_{ij}=_iA_j_jA_i`$ is the ordinary field strength (we are considering $`U(1)`$ case for simplicity) and $$\mathrm{\Psi }^i=\frac{1}{2}\left(\psi ^i+\overline{\psi }^i\right)$$ (1.211) The variation of (1.210) under the supersymmetry transformations (1.209) is a total derivative $$\delta 𝑑\tau \left(A_i(x)_\tau x^iiF_{ij}\mathrm{\Psi }^i\mathrm{\Psi }^j\right)=2i\eta 𝑑\tau _\tau (A_i\mathrm{\Psi }^i)$$ (1.212) Exactly as in bosonic case, the theory is regularized by making use of a “point splitting” technique and extra boundary terms are produced. Expanding the exponential of the action to first order in $`A`$ we can compute the variation of the path integral up to first order in $`L_A`$ $$i𝑑\tau 𝑑\tau ^{}\left(A_i_\tau x^i(\tau )iF_{ij}\mathrm{\Psi }^i\mathrm{\Psi }^j(\tau )\right)\left(2i\eta _\tau ^{}A_k\mathrm{\Psi }^k(\tau ^{})\right)$$ (1.213) Extra boundary terms appear when $`\tau ^{}\tau ^+`$ and $`\tau ^{}\tau ^{}`$. If the following interaction term is added to the action $$𝑑\tau A_iA_j\mathrm{\Psi }^i\mathrm{\Psi }^j(\tau )$$ (1.214) the extra terms are cancelled by its variation under (1.209). So we deduce that $`L_A`$ must be changed into $$i𝑑\tau \left(A_i(x)_\tau x^ii\widehat{F}_{ij}\mathrm{\Psi }^i\mathrm{\Psi }^j\right)$$ (1.215) being $`\widehat{F}`$ the noncommutative field strength. If we had performed a Pauli-Villars regularization, instead, (1.210) would have been invariant under worlsheet supersymmetry (1.209) and we would have ended up with field strength and gauge symmetry of ordinary $`U(1)`$. Summarizing, in it has been shown that the RNS open string in the presence of a constant $`B`$-field and D-branes can be described, in the zero-slope limit $`\alpha ^{}0`$, either by ordinary gauge theory or by noncommutative gauge theory on the brane, depending on the choice of the regularization prescription. The two different descriptions are related by a Seiberg-Witten map, as in bosonic case. In a different approach to the problem of coupling a gauge field to the open string in the presence of the $`B`$ field is presented. The coupling to the $`A`$ field is reconsidered in a way to preserve both shift symmetry and supersymmetry. A coupling term different with respect to (1.210) is found that is not supersymmetric by itself, but only together with the rest of the action $`S+S_b`$ and after making use of the corresponding boundary conditions. ##### Open GS superstring in the presence of a constant $`B`$-field and D-branes The Green-Schwarz superstring has manifest target space $`N=2`$ supersymmetry. The target space is a ten dimensional superspace described by the coordinates $`(x^\mu ,\theta ^{\alpha i})`$, with $`i=1,2`$. When the theory is coupled to a certain background and D-branes are present, in principle target space fermionic coordinates could be involved in noncommutativity. In , it has been shown that this is not true in the simple case of a constant NS-NS background. In this case only bosonic coordinates become noncommutative. In section 1.2.4 we will see that fermionic coordinates are indeed involved in noncommutativity when a constant R-R background is present. The action for the GS superstring coupled to a constant NS-NS background in flat spacetime is given by $`S_{GS}`$ $`={\displaystyle \frac{1}{2\pi \alpha ^{}}}{\displaystyle }d^2\xi \{\mathrm{\Pi }_i^\mu \mathrm{\Pi }^{i\nu }g_{\mu \nu }+2iϵ^{ij}_ix^\mu (\overline{\theta }^1\mathrm{\Gamma }_\mu _j\theta ^1\overline{\theta }^2\mathrm{\Gamma }_\mu _j\theta ^2)`$ (1.217) $`2ϵ^{ij}(\overline{\theta }^1\mathrm{\Gamma }^\mu _i\theta ^1)(\overline{\theta }^2\mathrm{\Gamma }_\mu _j\theta ^2)+ϵ^{ij}_ix^\mu _jx^\nu B_{\mu \nu }\}`$ where $$\mathrm{\Pi }_i^\mu =_ix^\mu i\overline{\theta }\mathrm{\Gamma }^\mu _i\theta $$ (1.218) are the supersymmetry invariant one-forms. It is clear from (1.217) that only bosonic coordinates couple to $`B_{\mu \nu }`$. However, boundary conditions along the D-brane must also be considered before we can deduce that fermions are not affected by the presence of the background. We find that bosonic coordinates must satisfy the same boundary conditions we found in the bosonic case (1.158). The two fermionic coordinates $`\theta ^1`$ and $`\theta ^2`$ must satisfy $`\theta ^2=\mathrm{\Gamma }_B\theta ^1`$ on the boundary, where $`\mathrm{\Gamma }_B`$ is a suitable $`B`$-dependent matrix satisfying $`\mathrm{\Gamma }_B^2=1`$. So the action (1.217) with these boundary conditions is invariant under the supersymmetry $`\delta x^\mu =i\overline{ϵ}\mathrm{\Gamma }^\mu \theta `$ (1.219) $`\delta \theta =ϵ`$ (1.220) only if the supersymmetry parameters $`ϵ^i`$ also satisfy $`ϵ^2=\mathrm{\Gamma }_Bϵ^1`$. So the D-brane breaks half of the supersymmetry. In an explicit computation is performed to dermine the coordinate algebra on the brane. Following the method introduced by Chu and Ho , the authors deal with boundary conditions along the brane by treating them as constraints on phase-space. The presence of these constraints makes it necessary to consider Dirac brackets instead of Poisson brackets. Working in a light-cone gauge, in it was shown that the fermionic variables surviving the gauge fixing procedure satisfy a standard anticommutative algebra and only the bosonic sector is affected by the presence of the constant NS-NS background, exactly as in . So it is clear that if we want target-space fermionic variables to be deformed, we must consider a different background. #### 1.2.3 Noncommutative selfdual Yang-Mills from the $`𝒩=2`$ string In this section I will briefly introduce the $`N=2`$ string and its peculiar properties. Referring to , I will apply the analysis outlined in the previous section to the open $`N=2`$ string in the presence of $`n`$ spacefilling D3-branes and a constant $`B`$ field to show that it coincides at tree level with $`U(n)`$ noncommutative selfdual Yang-Mills. The $`N=2`$ fermionic string is characterized by an $`N=2`$ worldsheet supersymmetry. For the string to be critical, the target space must be four-dimensional with signature $`(2,2)`$. Its propagating degrees of freedom are the embedding coordinates $`x^\mu `$ and the RNS Majorana spinors $`\psi ^\mu `$, $`\mu =1,2,3,4`$. These matter fields are coupled to the $`N=2`$ supergravity multiplet. Using symmetries of the action, one can show that all the gravitational degrees of freedom can be gauged away and in superconformal gauge the action can be written as follows $$S=\frac{1}{4\pi \alpha ^{}}_\mathrm{\Sigma }d^2\sigma \eta ^{\alpha \beta }\left(_\alpha x^\mu _\beta x^\nu +i\overline{\psi }^\mu \rho _\alpha _\beta \psi ^\nu \right)g_{\mu \nu }$$ (1.221) where $`g_{\mu \nu }=\zeta \mathrm{diag}(+1,+1,1,1)`$ ($`\zeta >0`$ scaling parameter) is the metric in $`R^{(2,2)}`$. This action has a residual symmetry given by the $`N=2`$ superconformal group, that also contains rigid $`N=2`$ supersymmetry corresponding to the transformations $`\delta x^\mu =\overline{ϵ}_1\psi ^\mu +J_\nu ^\mu \overline{ϵ}_2\psi ^\nu `$ (1.222) $`\delta \psi ^\mu =i\rho ^\alpha _\alpha x^\mu ϵ_1+iJ_\nu ^\mu \rho ^\alpha _\alpha x^\nu ϵ_2`$ (1.223) where $`J_\mu ^\nu `$ is a complex structure compatible with the metric $`g_{\mu \nu }J_\lambda ^\nu +J_\mu ^\nu g_{\lambda \nu }=0`$ (with our flat metric $`J_2^1=J_1^2=J_4^3=J_3^4=1`$). It has been shown that the open $`N=2`$ fermionic string at tree level is identical to self-dual Yang-Mills in $`2+2`$ dimensions . The absence of massive states in the physical spectrum is related to the vanishing of all tree-level amplitudes beyond three-point. The vanishing of amplitudes implies the existence of symmetries and vice-versa. Since an infinite number of tree-level amplitudes vanish, we expect an infinite number of symmetries to be present. These were described in . The $`N=2`$ string seems to be a master theory quantizing integrable systems. It is well-known that most integrable models in $`d=2`$ and $`d=3`$ can be obtained by dimensional reduction from selfdual Yang-Mills. In a possible definition for integrability in $`d=(2,2)`$ has been proposed, inspired by the peculiar properties of the tree-level S-matrix of the $`N=2`$ string. A system in $`d=(2,2)`$ would be classically integrable if the $`n`$-point tree-level amplitudes vanish beyond $`n=3`$. This definition is reminiscent of the one that can be given in $`d=2`$ concerning factorization of the S-matrix. The $`N=2`$ open fermionic string can be coupled to a two-form NS-NS background $`B`$ field. In it has been shown that additional boundary terms must be added to the action for the boundary conditions obtained from the Euler-Lagrange procedure to be $`N=2`$ supersymmetric, exactly as in the case of the $`N=1`$ string . Moreover, the presence of the second supersymmetry implies a nontrivial constraint on the $`B`$ field. This is the compatibility condition with respect to the complex structure $`B_{\mu \nu }J_\lambda ^\mu J_\mu ^\nu B_{\lambda \nu }=0`$, i.e. $`B`$ must be Kähler. The consistent $`N=2`$ gauge fixed action is then $`S`$ $`=`$ $`{\displaystyle \frac{1}{4\pi \alpha ^{}}}{\displaystyle _\mathrm{\Sigma }}d^2\sigma [(\eta ^{\alpha \beta }g_{\mu \nu }+ϵ^{\alpha \beta }2\pi \alpha ^{}B_{\mu \nu })_\alpha x^\mu _\beta x^\nu `$ (1.225) $`+(g_{\mu \nu }+2\pi \alpha ^{}B_{\mu \nu })i\overline{\psi }^\mu \rho _\alpha _\alpha \psi ^\nu ]`$ This action functional cannot be obtained from a superspace formulation. In it has been shown that the open $`N=2`$ string dynamics in the presence of $`n`$ coincident D3-branes filling the target space is modified by a magnetic $`B`$ field so that the open string sees an effective metric $`G_{\mu \nu }`$ and a noncommutative algebra on the brane characterized by a $`\theta ^{\mu \nu }`$ parameter. The starting point for the analysis is again the expression for the open string correlators (1.161). A particular choice of the $`SO(2,2)`$ generators allows us to write the matrices $`J`$ and $`B`$ in terms of the generators of the $`U(1)\times U(1)`$ subgroup of $`SO(2,2)`$, so that in complete generality the NS-NS field is expressed in terms of the two quantities $`B_1`$ and $`B_2`$ as $$B_{12}=B_{21}B_1B_{34}=B_{43}B_2$$ (1.226) In this basis the open string effective metric $`G^{\mu \nu }`$, noncommutativity parameter $`\theta ^{\mu \nu }`$ and coupling $`G_s`$ can be obtained and have expressions similar to the corresponding ones for the ten-dimensional string (1.164). Exactly as in the ten dimensional case a zero-slope limit can be taken that keeps the open string parameters finite. In this limit it has been shown that tree-level three-string amplitudes can be obtained from the noncommutative version of self-dual Yang-Mills in Leznov gauge (this depends on the choice of the Lorentz frame, one can as well obtain Yang gauge). The corresponding lagrangian is given by $`={\displaystyle \frac{1}{2}}G^{\mu \nu }\mathrm{tr}\left(_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }\right)+{\displaystyle \frac{1}{3}}ϵ^{\dot{\alpha }\dot{\beta }}\mathrm{tr}\left(\mathrm{\Phi }\widehat{}_{0\dot{\alpha }}\mathrm{\Phi }\widehat{}_{0\dot{\beta }}\mathrm{\Phi }\right)`$ (1.227) where $`G^{\mu \nu }e_{\widehat{\sigma }}^\mu e_{\widehat{\lambda }}^\nu \eta ^{\widehat{\sigma }\widehat{\lambda }}`$ is the open string effective metric and $`\widehat{}`$ is the corresponding derivative defined by $`\widehat{}_{\widehat{\mu }}e_{\widehat{\mu }}^\nu _\nu `$. This result is deduced as in the ten-dimensional case by noting that the effect of turning on the $`B`$-field is the multiplication of any open string amplitude by a phase factor (1.173). This corresponds to replacing ordinary products with Moyal products in the worldvolume effective field theory. As a nontrivial check, it has also been shown that the tree-level four-point function for $`U(n)`$ noncommutative self-dual Yang-Mills in Leznov gauge vanishes. Therefore, the natural deformation of selfdual Yang-Mills in Leznov gauge seems to preserve the nice scattering properties of the original, commutative theory. We have seen that the vanishing of tree-amplitudes beyond three-point defines integrable systems in $`d=(2,2)`$ exactly as the factorization of the S-matrix is a definition of a an integrable system in $`d=2`$. So the result in suggests that noncommutative selfdual Yang-Mills is integrable. This result will be important for the further developments considered in , where, in collaboration with O. Lechtenfeld, L. Mazzanti, S. Penati and A. Popov, I have constructed a noncommutative version of the sine-Gordon theory with a factorized S-matrix, that is obtained as a dimensional reduction of $`(2,2)`$ selfdual Yang Mills. #### 1.2.4 Non(anti)commutative field theories from the covariant superstring In this section I would like to discuss the superstring origin of the non(anti)commutative superspaces introduced in sections 1.1.5, 1.1.6, 1.1.7. Up to now I have discussed in detail the case of the open bosonic string in flat space in the presence of a constant NS-NS background field and Dp-branes and I have shown that the presence of the background induces a noncommutativity in the brane coordinate algebra. I have generalized this result to the RNS superstring, that is characterized by an $`N=1`$ worldsheet supersymmetry and happens to exhibit target space supersymmetry after a consistent truncation of the spectrum (GSO projection) is performed. This theory is not manifestly supersymmetric in target space. I have also shown how the bosonic results can be generalized to the manifestly target-space supersymmetric GS superstring. In this case we have seen that the presence of a constant NS-NS background does not modify the anticommutators between fermionic coordinates. Finally, I have generalized the bosonic string results to the $`N=2`$ string, characterized by an $`N=2`$ worldsheet supersymmetry. In all these cases a constant NS-NS two-form background has been considered. In this section we face the problem of finding a suitable superstring background that could induce, in the presence of D-branes, a nonanticommutative deformation of target space fermionic variables, exactly as the NS-NS B-field induces a deformation of bosonic target space coordinates. So target space must be a superspace and thus we need to consider a manifestly supersymmetric version of the superstring. An example of this is the Green-Schwarz superstring we considered in the previous section. It lives in a ten-dimensional superspace and its action is manifestly supersymmetric. Unfortunately, because of its complicated worldsheet symmetries, its action in a flat background is not quadratic, it cannot be quantized in a Lorentz-covariant way and this renders the formalism terribly difficult to handle. Recently a new proposal was made for an action describing the ten-dimensional superstring, that is manifestly target-space supersymmetric and is quadratic in a flat background<sup>4</sup><sup>4</sup>4An introduction to this formalism will be given in chapter 3.. In this formalism the R-R field strengths are all contained in a bispinor $`P^{\alpha \widehat{\alpha }}`$ where the two indices may have opposite or same chirality depending whether we are in IIA or IIB superstring theory. It is natural to think that, as much as the $`B^{\mu \nu }`$ background is related to a noncommutative deformation $`[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }`$, the $`P^{\alpha \widehat{\alpha }}`$ background may be related to a $`\{\theta ^\alpha ,\widehat{\theta }^{\widehat{\alpha }}\}=C^{\alpha \widehat{\alpha }}`$ deformation. In a series of papers it has been shown that this is indeed the case. In the first three papers the four-dimensional theory obtained from type II ten-dimensional superstring compactified on a Calabi-Yau three-fold has been considered, while in the last paper the full ten-dimensional superstring theory has been discussed. I will consider the four-dimensional case first, since, compared to the covariant quantization of the superstring in ten dimension, the formalism is much simpler in this case, because of the smaller amount of manifest supersymmetry. The relevant part of the lagrangian density is $$=\frac{1}{2}x^{\alpha \dot{\alpha }}\overline{}x_{\alpha \dot{\alpha }}+p_\alpha \overline{}\theta ^\alpha +p_{\dot{\alpha }}\overline{}\theta ^{\dot{\alpha }}+\overline{p}_\alpha \overline{\theta }^\alpha +\overline{p}_{\dot{\alpha }}\overline{\theta }^{\dot{\alpha }}$$ (1.228) where $`p_\alpha `$ $`\overline{p}_\alpha `$ $`p_{\dot{\alpha }}`$ and $`\overline{p}_{\dot{\alpha }}`$ are conjugate momenta to the superspace fermionic variables $`\theta ^\alpha `$ $`\overline{\theta }^\alpha `$ $`\theta ^{\dot{\alpha }}`$ and $`\overline{\theta }^{\dot{\alpha }}`$ respectively<sup>5</sup><sup>5</sup>5Berkovits covariant formalism is first order for fermionic variables, whose conjugate momenta are introduced as independent fields. This was first done by Siegel in his approach to the GS superstring . We indicate with dots target space Weyl spinor chirality and with bars worldsheet holomorphicity. The four dimensional action corresponding to (1.228) describes a free conformal field theory. The fields $`x`$, $`\theta `$, $`\overline{\theta }`$, $`p`$ and $`\overline{p}`$ satisfy free equations of motion, second order for $`x`$ and first order for fermionic variables. This theory exhibits an $`N=2`$ target-space supersymmetry<sup>6</sup><sup>6</sup>6It is interesting to note that, if one considers the undotted fermions only, the model can be regarded as a topological B-model on euclidean four-dimensional space and the topological BRST symmetry is strictly connected to the susy transformations generated by the dotted charges $`Q_{\dot{\alpha }}`$ and $`\overline{Q}_{\dot{\alpha }}`$.. It is useful to apply to (1.228) the following change of variables $`y^{\alpha \dot{\alpha }}=x^{\alpha \dot{\alpha }}+i\theta ^\alpha \theta ^{\dot{\alpha }}+i\overline{\theta }^\alpha \overline{\theta }^{\dot{\alpha }}`$ (1.229) $`q_\alpha =p_\alpha i\theta ^{\dot{\alpha }}x_{\alpha \dot{\alpha }}+{\displaystyle \frac{1}{2}}\theta ^{\dot{\beta }}\theta _{\dot{\beta }}\theta _\alpha {\displaystyle \frac{3}{2}}(\theta _\alpha \theta ^{\dot{\beta }}\theta _{\dot{\beta }})`$ (1.230) $`d_{\dot{\alpha }}=p_{\dot{\alpha }}i\theta ^\alpha x_{\alpha \dot{\alpha }}\theta ^\beta \theta _\beta \theta _{\dot{\alpha }}+{\displaystyle \frac{1}{2}}(\theta ^\beta \theta _\beta )`$ (1.231) and analogous for $`\overline{q}_\alpha `$ and $`\overline{d}_{\dot{\alpha }}`$. In the first line the reader will recognize the change of variables (1.148) introduced in section 1.1.7 to obtain Seiberg’s deformed superspace from my $`N=2`$ non(anti)commutative superspace. This transformation was first introduced by Vafa and Ooguri in . In the second and third lines I have written the worldsheet versions of the chiral supersymmetry charges and superspace covariant derivatives, $`Q_\alpha ={\displaystyle \frac{}{\theta ^\alpha }};\overline{Q}_\alpha ={\displaystyle \frac{}{\overline{\theta }^\alpha }}`$ (1.232) $`D_{\dot{\alpha }}={\displaystyle \frac{}{\theta ^{\dot{\alpha }}}};\overline{D}_{\dot{\alpha }}={\displaystyle \frac{}{\overline{\theta }^{\dot{\alpha }}}}`$ (1.233) ($`Q_{\dot{\alpha }}`$, $`\overline{Q}_{\dot{\alpha }}`$, $`D_\alpha `$ and $`\overline{D}_\alpha `$ are dressed in this representation). So, $`q`$ and $`d`$ represent the conjugate momenta to $`\theta `$’s at fixed $`y`$ exactly as $`p`$’s represent the same conjugate momenta at fixed $`x`$. The lagrangian (1.228) can be rewritten in terms of the new variables as follows $$=\frac{1}{2}y^{\alpha \dot{\alpha }}\overline{}y_{\alpha \dot{\alpha }}q_\alpha \overline{}\theta ^\alpha +d_{\dot{\alpha }}\overline{}\theta ^{\dot{\alpha }}\overline{q}_\alpha \overline{\theta }^\alpha +\overline{d}_{\dot{\alpha }}\overline{\theta }^{\dot{\alpha }}+\mathrm{total}\mathrm{derivative}$$ (1.234) In the presence of D-branes, one obtains the fermionic boundary conditions $`\theta ^\alpha =\overline{\theta }^\alpha ;q_\alpha =\overline{q}_\alpha `$ (1.235) $`\theta ^{\dot{\alpha }}=\overline{\theta }^{\dot{\alpha }};d_{\dot{\alpha }}=\overline{d}_{\dot{\alpha }}`$ (1.236) that only preserve half of the supersymmetry, generated by the charges $`Q_\alpha +\overline{Q}_\alpha `$ and $`Q_{\dot{\alpha }}+\overline{Q}_{\dot{\alpha }}`$. It is possible to couple the action to the R-R background described by the selfdual graviphoton field strength $$F^{\alpha \beta }0;F^{\dot{\alpha }\dot{\beta }}=0$$ (1.237) Only in euclidean signature it is possible to turn on the selfdual part of the super-two-form field strength $`F^{\alpha \beta }`$ while setting the antiselfdual part $`F^{\dot{\alpha }\dot{\beta }}`$ to zero. This is the stringy counterpart of the discussion we presented in section 1.1.5 for nonanticommutative superspaces and works exactly the same way. It is possible to show that the background (1.237) is an exact solution of the full nonlinear string equations of motion and that there is no backreaction to the metric. This can be seen from the fact that a purely selfdual field strength does not contribute to the energy momentum tensor and does not involve the dilaton field in its kinetic term. In the action the graviphoton field strength couples to the worldsheet supersymmetry currents $`q_\alpha `$ as follows $$F^{\alpha \beta }q_\alpha \overline{q}_\beta $$ (1.238) This is actually the graviphoton integrated vertex operator <sup>7</sup><sup>7</sup>7In chapter 3 I will give a detailed derivation of the corresponding vertex operator in the ten-dimensional case., that in this case gives the whole dynamics since there is no backreaction. We are interested in the effect of the background on the dynamics, so we can concentrate on the $`(q,\overline{q})`$ sector, with the lagrangian $$=\frac{1}{\alpha ^{}}\left(q_\alpha \overline{}\theta ^\alpha \overline{q}_\alpha \overline{\theta }^\alpha +\alpha ^{}F^{\alpha \beta }q_\alpha \overline{q}_\beta \right)$$ (1.239) We can integrate out the fields $`q_\alpha `$ and $`\overline{q}_\alpha `$ by using their equations of motion $`\overline{}\theta ^\alpha =\alpha ^{}F^{\alpha \beta }\overline{q}_\beta `$ (1.240) $`\overline{\theta }^\alpha =\alpha ^{}F^{\alpha \beta }q_\beta `$ (1.241) and obtain the effective lagrangian $$_{\mathrm{eff}}=\left(\frac{1}{\alpha _{}^{}{}_{}{}^{2}F}\right)_{\alpha \beta }\overline{\theta }^\alpha \overline{}\theta ^\beta $$ (1.242) We obtain boundary conditions for both fermionic variables and their derivatives $`\theta ^\alpha `$ $`=`$ $`\overline{\theta }^\alpha `$ (1.243) $`\overline{\theta }^\alpha `$ $`=`$ $`\overline{}\theta ^\alpha `$ (1.244) The first condition breaks half of the supersymmetry on the boundary. The second one corresponds to the equality of supersymmetry charges $`q_\alpha =\overline{q}_\alpha `$ on the boundary (see the equations of motion 1.241). The boundary conditions for derivatives of $`\theta `$ has first appeared in for the GS superstring. In that case they were additional, unexpected conditions, required by consistency of boundary conditions under kappa-symmetry. They were shown not to overconstrain the system, because they arise as restrictions of the field equations to the boundary. In this case instead the effective action (1.239), obtained by integrating out the fermionic conjugate momenta, is second order for the fermions. A boundary condition for derivatives of $`\theta `$ naturally arises from requiring that there are no surface terms in the Euler-Lagrange equations of motion. One can determine the fermionic propagators $`\theta ^\alpha (z)\theta ^\beta (w)={\displaystyle \frac{\alpha _{}^{}{}_{}{}^{2}F^{\alpha \beta }}{2\pi i}}\mathrm{log}{\displaystyle \frac{\overline{z}w}{z\overline{w}}}`$ (1.245) $`\overline{\theta }^\alpha (z)\overline{\theta }^\beta (w)={\displaystyle \frac{\alpha _{}^{}{}_{}{}^{2}F^{\alpha \beta }}{2\pi i}}\mathrm{log}{\displaystyle \frac{\overline{z}w}{z\overline{w}}}`$ (1.246) $`\theta ^\alpha (z)\overline{\theta }^\beta (w)={\displaystyle \frac{\alpha _{}^{}{}_{}{}^{2}F^{\alpha \beta }}{2\pi i}}\mathrm{log}{\displaystyle \frac{(zw)(\overline{z}\overline{w})}{(z\overline{w})^2}}`$ (1.247) On the boundary we get $$\theta ^\alpha (\tau )\theta ^\beta (\tau ^{})=\frac{\alpha _{}^{}{}_{}{}^{2}F^{\alpha \beta }}{2}ϵ(\tau \tau ^{})$$ (1.248) corresponding to the algebra $$\{\theta ^\alpha ,\theta ^\beta \}=\alpha _{}^{}{}_{}{}^{2}F^{\alpha \beta }$$ (1.249) Since the coordinates $`y`$ and $`\overline{\theta }`$ are not affected by the background coupling, they remain commuting. This means that the algebra written in terms of the original coordinate $`x`$ involves nontrivial terms in $`[x,x]`$ and $`[x,\theta ]`$, as required by consistency and first shown in my paper . In the deformation (1.249) was not welcome. It has been shown that when the open superstring is also coupled to a constant gluino superfield on the boundary, by adding the following term to the action $$W^\alpha q_\alpha $$ (1.250) the superspace deformation (1.249) is undone and supersymmetry is restored if the gluino fields satisfy the deformed algebra $$\{W_\alpha ,W_\beta \}=F^{\alpha \beta }$$ (1.251) Instead, in the supersymmetry-deforming algebra (1.249) was accepted and it was shown that in the zero-slope limit the $`N=\frac{1}{2}`$ theories we discussed in section 1.1.7 naturally emerge in string theory. The analysis pursued in was further developed in . In this paper it was shown that the constant selfdual background deforms the original $`N=2`$ superPoincaré algebra into another algebra that has still eight supercharges, four of which are unaffected by the background. In the presence of a D-brane, $`N=\frac{1}{2}`$ supersymmetry is realized linearly and the remaining $`N=\frac{3}{2}`$ is realized nonlinearly. This interpretation of the new terms arising in the supersymmetry algebra is similar to the original one we gave in . There we didn’t consider the new terms arising in the supersymmetry algebra as symmetry-breaking terms. We considered them as symmetry-deforming terms and we recast them in the known form of a q-deformation. It is very interesting to note that if both selfdual and antiselfdual field strength are turned on, there is a backreaction that warps spacetime to euclidean $`AdS_2\times S^2`$. The string in this background has been studied in and the structure of the action closely resembles the pure spinor version of the superstring in $`AdS_5\times S^5`$ . The action becomes quadratic in the limit $`F^{\dot{\alpha }\dot{\beta }}0`$ resembling the Penrose limit in the ten-dimensional case. $`N=\frac{1}{2}`$ super Yang-Mills on euclidean $`AdS_2\times S^2`$ has been studied in . A comparison to noncommutativity in the bosonic case is due. First of all, we notice that the superspace deformation (1.249) vanishes in the zero slope limit $`\alpha ^{}0`$ unless we also take the limit $`F^{\alpha \beta }\mathrm{}`$. This in principle can be done since $`F^{\alpha \beta }`$ is an exact solution to the string equations. Moreover, it is interesting to note that, in contradistinction to the bosonic case, where the $`B_{\mu \nu }`$ term only affects the boundary conditions, the $`F^{\alpha \beta }`$ term only affects the bulk equations of motion. It would be nice to see if there is a duality transformation connecting the two cases <sup>8</sup><sup>8</sup>8In S-duality was considered in the context of noncommutative geometry in the presence of both NS-NS B-field and R-R potentials.. The four-dimensional results in have been generalized to ten dimensions in . The main difference between the compactified four dimensional case and the full ten dimensional case is that a constant R-R field strength background, represented by a bispinor $`P^{\alpha \widehat{\alpha }}`$, is not a solution of the ten-dimensional string equations of motion (full nonlinear type II SUGRA equations), but it is only a solution of the linearized equations. Therefore in this case there is a backreaction. Nevertheless, it is possible to compute the corresponding vertex operator to be added to the action. The resulting theory is not the complete sigma model. In it was considered the general case where a NS-NS constant $`B^{\mu \nu }`$, a constant R-R field strength $`P^{\alpha \widehat{\alpha }}`$ and constant gravitinos $`\mathrm{\Psi }_m^\alpha `$ $`\mathrm{\Psi }_m^{\widehat{\alpha }}`$ are turned on. The resulting algebra is characterized by the usual nontrivial bosonic commutator induced by $`B`$, a nontrivial anticommutator between $`\theta ^\alpha `$ and $`\theta ^{\widehat{\alpha }}`$ induced by the R-R field strength and nontrivial commutators between the bosonic coordinate and the two fermionic ones induced by the two gravitinos. It would be nice to perform an analogous analysis in the pure spinor version of the superstring in a p-p wave background . In this case the whole sigma model action is known that involves a constant R-R field strength coupled to fermionic variables in a similar way with respect to (1.242). After the superspace deformation of a theory is known, one can integrate out the fermionic variables to obtain the component formulation of the theory. It has been shown in that the component formulation of $`N=\frac{1}{2}`$ super Yang-Mills theory can be obtained from standard RNS type IIB string theory compactified of a Calabi-Yau three-fold in the presence of a constant graviphoton background with a definite duality. A graviphoton background can be obtained in euclidean space by wrapping a R-R 5-form around a 3-cycle of internal Calabi-Yau space. Even if in the general case the RNS is not suitable to deal with a R-R background, in the constant graviphoton field strength case there are simplifications that allow for the computation of tree-level scattering amplitudes on the disk with the insertion of R-R vertex operators. The method used in is intrinsically perturbative, but the results are exact in the $`\alpha ^{}0`$ limit. So the nonanticommutative $`N=\frac{1}{2}`$ super Yang-Mills action in its component formulation is recovered from RNS string computations. ### 1.3 Generalization to non-constant backgrounds In this section I would like to discuss some results generalizing the connection between string theory and noncommutative geometry to the case of a nonconstant $`B`$ field. Let us introduce the notation $`\omega =F+B`$ for the gauge invariant combination in terms of which the physics of an open string in the presence of Dp-branes, gauge field $`A^i`$ and background $`B^{ij}`$ field is described. The latter will not be constrained to be constant from now on. We have to consider three possible situations, of growing complexity * constant $`\omega `$ * nonconstant $`\omega `$, $`d\omega H=0`$ * nonconstant $`\omega `$, $`d\omega H0`$ We have discussed the first case, corresponding to flat D-brane and flat background, in sections 1.2.1, 1.2.2, 1.2.3 . Summarizing the results in , in this case the Dp-brane worldvolume is described in terms of noncommuting coordinates. $`\omega `$ defines an associative symplectic structure associated to the noncommutativity matrix $`\theta `$. In the limit $`\alpha ^{}0`$ physics can be described by an effective theory which is a gauge theory deformed by the noncommutative associative Moyal product. In the limit $`\alpha ^{}0`$ the noncommutative parameter $`\theta `$ is given by the inverse of $`\omega `$ (see 1.180). The second case has been studied in . It describes a physical situation with curved Dp-branes in a flat background. The worldvolume deformation is decribed by Kontsevich $``$ product, where the noncommutativity parameter is given by the inverse of $`\omega `$, as before. In this case $`\omega `$ is not constant anymore, so neither is $`\theta `$ and the correct product is Kontsevich $``$. $`\omega ^1`$ is still a Poisson structure on the manifold, that now exhibits a coordinate dependence such that associativity is satisfied. The formula (1.185), valid in the limit $`\alpha ^{}0`$, is generalized to $$\underset{i=1}{\overset{n}{}}f_i\left(x(\tau _i)\right)=V(\omega )𝑑xf_1\mathrm{}f_n$$ (1.252) being $``$ Kontsevich associative product. The last case, discussed in , is a further generalization to the case where the background is also curved. $`\omega `$ is not a Poisson structure anymore, since associativity is lost. This is related to the emergence of a second geometrical object playing a role in the physics of the system. This is the 3-form $`H=d\omega `$, that has been shown to be the parameter governing nonassociativity. However, it is always possible to give a description of the Dp-brane worldvolume in terms of a Kontsevich-like product. Noncommutativity is still governed by $`\omega ^1`$ and nonassociativity by $`H`$. In a metric $`g^{\mu \nu }`$ is considered that is a small perturbation from the flat metric of section 1.2.1. The string action in a generic curved background is still of the form (1.157), where the target space metric is $`g^{\mu \nu }(x)`$, describing a curved spacetime. It is possible to expand the action around the flat spacetime metric as follows $$S=S_0+S_1+\mathrm{}$$ (1.253) where $`S_0`$ represents the action in a flat spacetime (1.157) and $`S_1`$ is the first correction due to the presence of a small curvature. In situations where a curved background is present, it is convenient to choose special coordinates, known as Riemann normal coordinates, defined along geodesics in target space starting from $`x^\mu =0`$ (see for instance , where the normal coordinate expansion of the metric is introduced in the context of sigma models and string theory). In terms of these coordinates the Taylor expansion of every tensor around $`x^\mu =0`$ is expressed in terms of covariant tensors evaluated at the origin. In particular, the expansion for the metric up to second order in $`x`$ is given by $$g_{\mu \nu }(x)=g_{\mu \nu }\frac{1}{3}R_{\mu \rho \nu \sigma }x^\rho x^\sigma +𝒪(x^3)$$ (1.254) where $`R_{\mu \nu \rho \sigma }`$ is the curvature tensor. The analogous expansion for the $`B`$ field, in radial gauge, is $$B_{ij}(x)=B_{ij}+\frac{1}{3}H_{ijk}x^k+\frac{1}{4}_lH_{ijk}x^kx^l+𝒪(x^3)$$ (1.255) with $`H=dB`$. In a first order approximation in $`x`$ is applied. The action considered is $$S=S_0+S_1+S_B$$ (1.256) where $`S_0+S_B`$ is the action (1.157) for flat $`g`$ and constant $`B`$, while $`S_1`$ is the first order correction in $`x`$, deriving from the $`B`$ expansion (1.255) $$S_1=\frac{i}{6}_\mathrm{\Sigma }H_{ijk}x^kϵ^{ab}_ax^i_bx^j$$ (1.257) $`S_1`$ is treated as an interaction term with respect to the free theory, described by the action $`S_0`$. The perturbative analysis with the interaction $`S_1`$ gives the following results. In the limit $`\alpha ^{}0`$ correlators are expressed in a form analogous to (1.252) $$\underset{i=1}{\overset{n}{}}f_i\left(x(\tau _i)\right)=V(\omega )d^{p+1}x(f_1\mathrm{}f_n)$$ (1.258) where $``$ is Kontsevich-like nonassociative product, whose definition is completely analogous to (1.78, 1.79), with $`P(x)=\omega ^1(x)`$. The relation (1.70) is not valid, though, for the noncommutativity parameter $`P=\omega ^1`$ and the violation of associativity is proportional to the 3-form $`H`$ as follows $$(fg)hf(gh)=\frac{1}{6}P^{im}P^{jn}P^{kl}H_{mnl}_if_jg_kh+\mathrm{}$$ (1.259) Note that (1.258) is problematic, because of the nonassociativity of $``$. It is necessary to specify the positions of the points $`\tau _i`$ on the boundary of the disk, so invariance under cyclic permutations is lost. The coordinate algebra is given by $$[x^i,x^j]_{}=iP^{ij}(x)$$ (1.260) showing that the role played by the two-form $`\omega `$ is unchanged in the general curved case. The relations concerning the nonassociativity of the coordinate algebra are instead $$\left(x^ix^j\right)x^kx^i\left(x^jx^k\right)=\frac{1}{6}P^{im}P^{jn}P^{kl}H_{mnl}$$ (1.261) $$[x^i,[x^j,x^k]_{}]_{}+[x^k,[x^i,x^j]_{}]_{}+[x^j,[x^k,x^i]_{}]_{}=P^{im}P^{jn}P^{kl}H_{mnl}$$ (1.262) These formulas further clarify the role played by the 3-form $`H`$ as the parameter governing nonassociativity. The deep relationship between spacetime geometry and nonassociativity on the D-brane worldvolume discovered in is very interesting, showing that nonassociative geometry also plays a role in the new developments regarding spacetime pioneered by string theory. It would be nice to generalize the results discussed in this section to the manifestly target space supersymmetric string (in Green-Schwarz or Berkovits formalism) to see whether a Kontsevich-like product similar to the one I proposed in , in collaboration with D. Klemm and S. Penati, may emerge in the presence of a nontrivial super-three-form field strength background . It would be also interesting to consider a generalization of the discussion presented in section 1.2.4, where the R-R field strength is not constant. In particular, when the R-R field strength has a linear dependence on the bosonic coordinates, it would be nice to investigate whether Lie algebraic deformed superspaces, characterized by the fermionic anticommutators $$\{\theta ^\alpha ,\theta ^\beta \}=\gamma _\mu ^{\alpha \beta }x^\mu $$ (1.263) can appear. These kind of superspace was first studied in , in a different context. There, the authors considered the possibility that bosonic spacetime had a fermionic substructure, given by the relation (1.263). In , in collaboration with P. A. Grassi, I have started to consider this problem by computing the vertex operator for the Berkovits covariant superstring for a R-R field strength with a linear dependence on the bosonic coordinates. This will be discussed in detail in chapter 3. ## Chapter 2 Noncommutative deformation of integrable field theories ### 2.1 A brief introduction to selected topics concerning two-dimensional classical integrable systems Disclaimer: This section is not a thorough introduction to the vast field of integrable systems. I will only review topics needed in the two following sections, where my study of the noncommutative sine-Gordon system will be presented. In most cases I will not give details and I will just give references for the interested reader. Moreover, for any single topic in this section I will exhibit a single example, the ordinary sine-Gordon model. #### 2.1.1 Infinite conserved currents and the bicomplex approach In classical mechanics a system described by $`n`$ degrees of freedom is completely integrable when it is endowed with $`n`$ conserved currents. In classical field theory, a system with an infinite number of local conserved currents is also said to be integrable. This is a property of the equations of motion. For some integrable system an action is also known that generates the equations by an Euler-Lagrange procedure, but this is not true in general. Indeed, it is true for the sine-Gordon model that I will consider in the rest of this chapter. By making use of the bicomplex technique it is possible to construct second order differential equations that are integrable. Moreover, with this approach it is very easy to generate the corresponding conserved currents by an iterative procedure. In this section we use euclidean signature and complex coordinates $$z=\frac{1}{\sqrt{2}}(x^0+ix^1);\overline{z}=\frac{1}{\sqrt{2}}(x^0ix^1)$$ (2.1) A bicomplex is a triple $`(,d,\delta )`$ where $`=_{r0}^r`$ is an $`N_0`$-graded associative (but not necessarily commutative) algebra, $`^0`$ is the algebra of functions on $`^2`$ and $`d,\delta :^r^{r+1}`$ are two linear maps satisfying the conditions $`d^2=\delta ^2=\{d,\delta \}=0`$. $`^r`$ is therefore a space of $`r`$-forms. Let us consider the linear equation $$\delta \xi =ld\xi $$ (2.2) where $`l`$ is a real parameter and $`\xi ^s`$ for a given “spin” $`s`$. Suppose a nontrivial solution $`\stackrel{~}{\xi }`$ exists. Expanding it in powers of the given parameter $`l`$ as $$\stackrel{~}{\xi }=\underset{i=0}{\overset{\mathrm{}}{}}l^i\xi ^{(i)}$$ (2.3) one obtains the following equations relating the components $`\xi ^{(i)}^s`$ $$\delta \xi ^{(0)}=0;\delta \xi ^{(i)}=d\xi ^{(i1)},i1$$ (2.4) Therefore we obtain the chain of $`\delta `$-closed and $`\delta `$-exact forms $$\mathrm{\Xi }^{(i+1)}d\xi ^{(i)}=\delta \xi ^{(i+1)},i0$$ (2.5) For the chain not to be trivial $`\xi ^{(0)}`$ must not be $`\delta `$-exact. When the two differential maps $`d`$ and $`\delta `$ are defined in terms of ordinary derivatives in $`^2`$, the conditions $`d^2=\delta ^2=\{d,\delta \}=0`$ are trivially satisfied. Therefore, the possibly infinite set of conservation laws (2.4) is not associated to any second order differential equation and it is not useful for our purpose. However, it is possible to gauge the bicomplex by dressing the two differential maps $`d`$ and $`\delta `$ with the connections $`A`$ and $`B`$ as follows $$D_d=d+A;D_\delta =\delta +B$$ (2.6) The flatness conditions $`D_d^2=D_\delta ^2=\{D_d,D_\delta \}=0`$ are now nontrivial and give the differential equations $`(A)dA+A^2=0`$ (2.7) $`(B)\delta B+B^2=0`$ (2.8) $`𝒢(A,B)dB+\delta A+\{A,B\}=0`$ (2.9) Exactly as before we can consider the linear equation corresponding to (2.2) $$𝒟\xi (D_\delta lD_d)\xi =0$$ (2.10) The nonlinear equations (2.9) are the compatibility conditions for (2.10) $$0=𝒟^2\xi =\left[(B)+l^2(A)l𝒢(A,B)\right]\xi $$ (2.11) Supposing that the linear equation (2.10) admits a solution $`\stackrel{~}{\xi }^s`$ and expanding it as $`\stackrel{~}{\xi }=_{i=0}^{\mathrm{}}l^i\xi ^{(i)}`$, one obtains the possibly infinite chain of identities $$D_\delta \xi ^{(0)}=0;D_\delta \xi ^{(i)}=D_d\xi ^{(i1)},i1$$ (2.12) from which $`D_\delta `$-closed and $`D_\delta `$-exact forms $`\mathrm{\Xi }^{(i)}`$ can be constructed, when $`\xi ^{(0)}`$ is not $`\delta `$-exact. Therefore, with a suitable choice of $`A`$ and $`B`$, it is possible to construct interesting second order integrable differential equations from (2.9) and their conserved currents from (2.12). In general the currents $`\xi ^{(i)}`$are nonlocal functions of the coordinates, since they may be expressed in terms of integrals, but it is possible to extract local currents from them that have a physical meaning. As an example, we can derive the ordinary sine-Gordon equation from this formalism. Let $`=_0L`$, where $`_0`$ is the space of $`2\times 2`$ matrices with entries in the algebra of smooth functions on ordinary $`^2`$ and $`L=_{i=0}^2L^i`$ is a two-dimensional graded vector space with the $`L^1`$ basis $`(\tau ,\sigma )`$ satisfying $`\tau ^2=\sigma ^2=\{\tau ,\sigma \}=0`$. If we take the differential maps $$\delta f=\overline{}f\tau Rf\sigma ;df=Sf\tau +f\sigma $$ (2.13) with commuting constant matrices $`R`$ and $`S`$, then the conditions $`d^2=\delta ^2=\{d,\delta \}=0`$ are trivially satisfied. To obtain nontrivial second order differential equations we can gauge the bicomplex by dressing $`d`$ as follows $$D_dfG^1d(Gf)$$ (2.14) with $`G`$ generic invertible matrix in $`_0`$. The condition $`D_d^2=0`$ is trivially satisfied, while $`\{\delta ,D_d\}=0`$ yields the nontrivial second order differential equation $$\overline{}\left(G^1G\right)=[R,G^1SG]$$ (2.15) With the choice $$R=S=2\sqrt{\gamma }\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)$$ (2.16) and taking $`GSU(2)`$ as follows $$G=e^{\frac{\mathrm{i}}{2}\sigma _2\mathrm{\Phi }}=\left(\begin{array}{cc}\mathrm{cos}\frac{\mathrm{\Phi }}{2}& \mathrm{sin}\frac{\mathrm{\Phi }}{2}\\ \mathrm{sin}\frac{\mathrm{\Phi }}{2}& \mathrm{cos}\frac{\mathrm{\Phi }}{2}\end{array}\right)$$ (2.17) we obtain the sine-Gordon equation from the off-diagonal part of the matrix equation (2.14) $$\overline{}\mathrm{\Phi }=4\gamma \mathrm{sin}\mathrm{\Phi }=0$$ (2.18) The diagonal part instead gives an equation which is trivially satisfied. From this derivation it is clear that the bicomplex approach can be naturally extended to noncommutative space, by replacing ordinary products with Moyal products in the whole discussion. In particular $`D_d=d+A`$ and $`D_\delta =\delta +B`$ and equations (2.9) are generalized accordingly. This will be discussed for the sine-Gordon model in section 2.2.4. Since in the noncommutative case the group $`SU(2)`$ is not closed any more and must be extended to $`U(2)`$, as we have seen in section 1.1.2, it is natural to expect that the noncommutative generalization of this construction for the sine-Gordon equation will be nontrivial. #### 2.1.2 Reductions from selfdual Yang-Mills The self-duality equations for Yang-Mills fields in $`^4`$ with signature $`(++++)`$ or $`(++)`$ $`{\displaystyle \frac{1}{2}}ϵ_{\mu \nu \rho \sigma }F^{\rho \sigma }=F_{\mu \nu }`$ (2.19) $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu +[A_\mu ,A_\nu ]`$ (2.20) are a famous example of nonlinear integrable equations in four dimensions. For $`SU(n)`$ gauge theory the potentials $`A_\mu ^a`$ are real. It is possible to consider an analytic continuation of $`A_\mu ^a`$ into complex space parametrized by the complex coordinates $`y`$, $`\overline{y}`$, $`z`$, $`\overline{z}`$. The selfduality equations (2.20) can be rewritten in the form $$F_{yz}=F_{\overline{y}\overline{z}}=0;F_{y\overline{y}}\pm F_{z\overline{z}}=0$$ (2.21) where the last sign is $`+`$ is the euclidean case and $``$ in the kleinian case. The first equation in (2.21) implies that the potentials $`A_y`$, $`A_z`$ ($`A_{\overline{y}}`$, $`A_{\overline{z}}`$) are pure gauges for fixed $`\overline{y}`$, $`\overline{z}`$ ($`y`$, $`z`$), therefore two $`n\times n`$ complex matrices $`B`$ and $`\overline{B}`$ can be found such that $`A_y=B^1_yB;A_z=B^1_zB`$ (2.22) $`A_{\overline{y}}=\overline{B}^1_{\overline{y}}\overline{B};A_{\overline{z}}=\overline{B}^1_{\overline{z}}\overline{B}`$ (2.23) Defining $`J=B\overline{B}^1SL(n,)`$, the last equation can be rewritten as $$_{\overline{y}}(J^1_yJ)\pm _{\overline{z}}(J^1_zJ)=0$$ (2.24) describing selfdual Yang-Mills in Yang formulation. This equation resembles the sum of two WZW model equations (see section 1.1.3) involving $`(y,\overline{y})`$ and $`(z,\overline{z})`$ variables, respectively. Therefore Yang equation (2.24) can be obtained from the following action $`S`$ $`=`$ $`{\displaystyle d^2yd^2z\mathrm{tr}(_yJ_{\overline{y}}J^1)}{\displaystyle d^2yd^2z_0^1𝑑\rho \mathrm{tr}\left(\widehat{J}^1_\rho \widehat{J}[\widehat{J}^1_{\overline{y}}\widehat{J},\widehat{J}^1_y\widehat{J}]\right)}`$ (2.25) $`+`$ $`{\displaystyle d^2yd^2z\mathrm{tr}(_zJ_{\overline{z}}J^1)}{\displaystyle d^2yd^2z_0^1𝑑\rho \mathrm{tr}\left(\widehat{J}^1_\rho \widehat{J}[\widehat{J}^1_{\overline{z}}\widehat{J},\widehat{J}^1_z\widehat{J}]\right)}`$ (2.27) where $`\widehat{J}(y,\overline{y},z,\overline{z},\rho )`$ is a homotopy path satisfying $`\widehat{J}(\rho =0)=1`$ and $`\widehat{J}(\rho =1)=J`$. Ward conjectured that all integrable equations in $`d=2`$ can be obtained as dimensional reductions of selfdual Yang-Mills equations , so that the latter play the role of a universal integrable system. Since then the conjecture has been tested on many integrable systems whose Lax pair can also be obtained by reduction from the one associated to the selfduality equations (2.20). The kleinian case is particularly interesting because of its connections with the $`N=2`$ string, discussed in section 1.2.2. Reductions are obtained by requiring invariance under any arbitrary subgroup $`G`$ of the group of conformal transformations of $`^{(2,2)}`$ (or $`^{(4,0)}`$). Afterwards algebraic constraints can be applied to the arbitrary matrices involved in the equations to obtain known integrable models. In most cases invariance under translations in certain directions is required. Therefore it is clear that there are many more possibilities in reducing selfdual Yang-Mills in $`^{(2,2)}`$ with respect to the euclidean case, since, instead of requiring invariance under the usual complex coordinates $`\sqrt{2}y=x^1+ix^2;\sqrt{2}\overline{y}=x^1ix^2`$ (2.28) $`\sqrt{2}z=x^3+ix^4;\sqrt{2}\overline{z}=x^3ix^4`$ (2.29) (combining space with space and time with time), in the kleinian case one can also require invariance with respect to light-cone coordinates $`s={\displaystyle \frac{1}{2}}(x^2x^4);t={\displaystyle \frac{1}{2}}(x^2+x^4)`$ (2.30) $`u={\displaystyle \frac{1}{2}}(x^1x^3);v={\displaystyle \frac{1}{2}}(x^1+x^3)`$ (2.31) (combining space and time). As an example I will show how ordinary euclidean sine-Gordon model can be obtained from selfdual Yang-Mills equations. The euclidean version of Yang equation (2.24) gives the sine-Gordon equation if one choses $$B=e^{\frac{z}{2}\sigma _1}e^{i\frac{\mathrm{\Phi }}{2}\sigma _3};\overline{B}=e^{\frac{\overline{z}}{2}\sigma _1}$$ (2.32) where $`\mathrm{\Phi }=\mathrm{\Phi }(y,\overline{y})`$. In fact one finds that the field $`\mathrm{\Phi }`$ satisfies (2.18) with $`4\gamma =1`$. It will be also useful to know that the sine-Gordon equation can be obtained from kleinian selfdual Yang-Mills equations through a two-step reduction procedure. First of all one requires independence of Yang equation (2.24) under one of the real coordinates $`x^i`$, reducing to the $`2+1`$ model $$(\eta ^{\mu \nu }+V_\alpha ϵ^{\alpha \mu \nu })_\mu (J^1_\nu J)=0$$ (2.33) where $`V_\alpha `$ is a constant vector in spacetime. A nonzero $`V_\alpha `$ breaks Lorentz invariance but restores integrability when it is spacelike and with unit length (nonlinear sigma models in $`2+1`$ dimensions can be Lorentz invariant or integrable but not both ). From this equation, in the case $`V_\alpha =(0,1,0)`$, we can make a further reduction by choosing $$J=\left(\begin{array}{cc}\mathrm{cos}\frac{\mathrm{\Phi }}{2}& e^{\frac{i}{2}x}\mathrm{sin}\frac{\mathrm{\Phi }}{2}\\ e^{\frac{i}{2}x}\mathrm{sin}\frac{\mathrm{\Phi }}{2}& \mathrm{cos}\frac{\mathrm{\Phi }}{2}\end{array}\right)SU(2)$$ (2.34) with $`\mathrm{\Phi }`$ depending on only two coordinates with different signature. As a result, we obtain the sine Gordon equation in $`1+1`$ dimensions for the field $`\mathrm{\Phi }`$. #### 2.1.3 Properties of the S-matrix In section 2.1.1 we have seen how to construct nonlinear differential equations in two-dimensions with the property of having an infinite number of conserved currents. We have also said that from that construction it is possible to extract conserved currents that are local and thus have a physical meaning. In this case the corresponding equations are called integrable. Moreover, we have seen an example of two-dimensional integrable system, the ordinary sine-Gordon equation. The S-matrix of a two-dimensional theory with an infinite set of conserved currents that are local and yield conserved charges which are components of Lorentz tensors of increasing rank enjoys several nice properties. * The general multiparticle S-matrix is elastic, i.e. the number and set of momenta of particles of any given mass remains the same before and after the collision. * Any multiparticle S-matrix factorizes into a product of two-particle S-matrices. * These two-particle S-matrices satisfy a cubic equation that in most cases is sufficient to obtain exact expressions for them (unitarity must be used, though!). The sine-Gordon model is one of these cases. These are of course very nice properties. Indeed, the task of computing the general S-matrix considerably simplifies, since it reduces to determining only the two-body S-matrix. I’m not going to give a proof of this theorem. The interested reader can for instance refer to . The key ingredient in the proof is the sufficient complexity of the conserved charges (i.e. the growing Lorentz rank). However, locality is strongly used and its is unclear whether the kind of nonlocality introduced by the $``$ product can be a problem. As we discussed in section 1.1.2, unitarity and causality problems arise in theories with noncommutating time, such as 1+1 theories, together with a general breakdown of the quantum mechanics framework. In my paper , an explicit example of a noncommutative two dimensional theory with an infinite number of conserved currents that does not have a factorized S-matrix was constructed. Therefore, it seems that the theorem cannot be trivially extended to noncommutative case. However, in my paper , a noncommutative two-dimensional system which is integrable and has a factorized S-matrix was constructed. It might be that a nontrivial interplay between integrability and causality drives a system to exhibit or not a factorized S-matrix. These issues will be discussed in the rest of this chapter. #### 2.1.4 Solitons The name soliton refers to solutions of nonlinear equations that represent a localized packet travelling without changing shape or velocity and preserving these properties after collision with other packets. In complicated equations containing nonlinear and dispersive terms, the existence of this kind of solution is a highly nontrivial property, due to a special balance between the effects of these two kinds of contributions. A very famous example of a system displaying this kind of classical solutions is the sine-Gordon model. It can be shown that this system, described by the equation of motion (2.18), admits the following static finite-energy solution $$\mathrm{\Phi }\mathrm{tan}^1[\mathrm{exp}(xx_0)]=\mathrm{\Phi }_{\mathrm{sol}}(xx_0)$$ (2.35) (soliton) and the corresponding one (antisoliton) obtained by the discrete transformation $`\mathrm{\Phi }\mathrm{\Phi }`$, which is a symmetry of the system. Moving solutions can be obtained from static ones by Lorentz transformation. A third kind of solution is present, called doublet or breather solution, which can be interpreted as a bound system made of a soliton-antisoliton pair. For a detailed derivation of these solutions the reader should refer to . It can be shown by studying exact time-dependent solutions representing scattering of solitons that the colliding solitons do not change shape or velocity after collision. From direct inspection of these scattering solutions, representing two (or more) (anti)solitons far apart and approaching with a relative velocity, it is clear that the only effect of the collision in the distant future is some time delay (see ). Solitons solutions are not present in any scalar field theory with a potential bounded from below in spatial dimension greater than two, as the energy of any field configuration can always be lowered by shrinking. This follows from a simple scaling argument by Derrick . Let us consider a theory for a single scalar field in D+1 dimensions for simplicity (the discussion can be easily extended to a set of $`N`$ interacting scalar fields), described by the standard relativistic lagrangian. The corresponding energy functional for static configurations is $$E=\frac{1}{g^2}d^Dx\left(\frac{1}{2}(\mathrm{\Phi })^2+V(\mathrm{\Phi })\right)$$ (2.36) Let $`\mathrm{\Phi }_0(x)`$ be an extremum of (2.36). Consider the energy of the configuration $`\mathrm{\Phi }_\lambda (x)=\mathrm{\Phi }_0(\lambda x)`$ $$E(\lambda )=\frac{1}{g^2}d^Dx\left(\frac{1}{2}\lambda ^{2D}(\mathrm{\Phi }_0(x))^2+\lambda ^DV(\mathrm{\Phi }_0(x))\right)$$ (2.37) Since we assumed that $`\mathrm{\Phi }_0(x)`$ is an extremum, we require $`\frac{E(\lambda )}{\lambda }\|_{\lambda =1}=0`$. This gives the equation $$\frac{1}{g^2}d^Dx\left(\frac{1}{2}(D2)(\mathrm{\Phi }_0(x))^2+DV(\mathrm{\Phi }_0(x))\right)=0$$ (2.38) If the potential $`V`$ is bounded from below by zero and $`D3`$, than kinetic and potential terms in this equation must vanish separately and thus no nontrivial space-dependent solutions are admitted. In the case $`D=2`$ one obtains that $`V(\mathrm{\Phi }_0(x))=0`$. If $`V`$ has only discrete minima, then also in $`D=2`$ no time-dependent solutions are allowed. However, when $`V`$ has a continuos set of minima (in the case with more than one scalar field), possible space-dependent solutions are permitted. Notice that this proof is not valid when higher derivative terms are present and when the scalar field is described by a nonrelativistic lagrangian. Moreover, only static solutions are excluded, while time-dependent ones are allowed. ### 2.2 Deforming an integrable field theory: The sine-Gordon model (a first attempt) #### 2.2.1 Noncommutative solitons In this section I will mostly refer to the ICTP lectures by R. Gopakumar . To the interested reader, I also suggest the lectures by N. Nekrasov and by J. Harvey . Since the literature concerning noncommutative solitons is vast, I only give a partial list of references in . To this, one should add references where solitons of a specific 2+1 integrable model are studied, which are related the noncommutative sine-Gordon solitons in my work that will be discussed in section 2.3.5. A quite universal feature of noncommutative field theories is that they admit classical finite energy soliton solutions that have no counterpart in local field theories. This novel soliton solutions are more or less insensitive to the details of the specific theory considered, so in this section I will consider the scalar theory in 2+1 dimensions for simplicity, with only spatial noncommutativity. In section 2.1.5 we have seen that ordinary scalar theory, with a standard relativistic lagrangian and a potential with a discrete set of minima, does not have any localized solution in spatial dimension greater than one (see section 2.1.5). In the following we will see that spatial nonlocality induced by noncommutativity allows for the presence of novel localized solutions that vanish in the commutative limit. Consider the energy functional for static configurations $$E=\frac{1}{g^2}d^2z\left(\mathrm{\Phi }\overline{}\mathrm{\Phi }+V(\mathrm{\Phi })_{}\right)$$ (2.39) where $`z`$, $`\overline{z}`$ are complex coordinates in the two dimensional noncommutative space and $``$ is the corresponding Moyal product. As we have seen in section 1.1.1, integrated quadratic terms are unaffected by Moyal product, so only the potential term is modified with respect to the ordinary theory. Since we know that for $`\theta =0`$ there are no solitonic solutions, we will first consider the limit $`\theta \mathrm{}`$. It is useful to rescale the complex coordinates $`zz\sqrt{\theta }`$, $`\overline{z}\sqrt{\theta }`$, so that the $``$ product does not depend explicitly on the deformation parameter and the energy functional becomes $$E=\frac{1}{g^2}d^2z\left(\mathrm{\Phi }\overline{}\mathrm{\Phi }+\theta V(\mathrm{\Phi })_{}\right)$$ (2.40) where all the $`\theta `$ dependence is in front of the potential term. In the limit $`\theta \mathrm{}`$ the kinetic term is negligible with respect to the potential term, at least for localized configurations varying over a size of order one in the rescaled coordinates. Therefore, we will look for solutions of $$\left(\frac{V}{\mathrm{\Phi }}\right)_{}=0$$ (2.41) For instance, in the case of a cubic potential, one has to solve the equation $$m^2\mathrm{\Phi }+b_3\mathrm{\Phi }\mathrm{\Phi }=0$$ (2.42) In the commutative case this equation would only admit the constant solutions $$\mathrm{\Phi }=\lambda _i$$ (2.43) where $`\lambda _i`$ are the extrema of the function $`V(\mathrm{\Phi })`$. Nonlocality introduced by Moyal product allows for more interesting solutions. Recalling the Weyl-Moyal correspondence we introduced in section 1.1.1, relating functions in a noncommutative algebra to operators in a suitable Hilbert space, we see that functions $`\mathrm{\Phi }(z,\overline{z})`$ satisfying $`\mathrm{\Phi }\mathrm{\Phi }=\mathrm{\Phi }`$ exist and correspond to projectors $`P`$, $`P^2=P`$, in the Hilbert space. Therefore it is clear that $`\widehat{\mathrm{\Phi }}=\lambda _iP`$ is a solution of (2.42) when $`P`$ is a projection operator on some subspace of the Hilbert space and $`\lambda _i`$ is an extremum of $`V`$. Since integration over coordinates $`z`$, $`\overline{z}`$ corresponds through the Weyl-Moyal correspondence to trace over the Hilbert space, the energy functional in operator language is $$E=\frac{2\pi \theta }{g^2}V(\lambda _i)\mathrm{trP}$$ (2.44) The most general solution to (2.42) is $$\widehat{\mathrm{\Phi }}=\underset{k}{}a_kP_k$$ (2.45) where the coefficients $`a_k`$ are chosen among the ordinary constant extrema $`\lambda _i`$ of $`V(\mathrm{\Phi })`$ and $`P_k`$ are mutually orthogonal projection operators. To understand the physical meaning of the solutions we have found, we have to go back to configuration space. One finds that the solutions (2.45) are radially symmetric in space and with an $`r`$-dependence given by $$\underset{n=0}{\overset{\mathrm{}}{}}a_n\varphi _n(r^2)$$ (2.46) where $$\varphi _n(r^2)=2(1)^ne^{r^2}L_n(2r^2)$$ (2.47) and $`L_n(x)`$ is the n-th Laguerre polynomial. The simplest solution $`\varphi _0(r)`$ is a gaussian. Moreover, non radially symmetric solutions can be generated by noting that the action in the limit $`\theta \mathrm{}`$ has the $`U(\mathrm{})`$ symmetry $`\widehat{\mathrm{\Phi }}U\mathrm{\Phi }U^{}`$, where $`U`$ is a unitary operator acting on the Hilbert space. It can be easily proven that radially symmetric solutions are stable against small fluctuations when they are constructed around a local minimum configuration $`\lambda _i`$ of the potential $`V`$ and that non radially symmetric are stable too, since $`U(\mathrm{})`$ rotations do not change the energy. So stable solitons are present when the potential has at least two minima. The $`U(\mathrm{})`$ symmetry is broken when $`\frac{1}{\theta }`$ corrections are taken into account (i.e. the kinetic term is not negligible anymore). Most of the infinite solutions we found in the $`\theta =\mathrm{}`$ case disappear, but it was found that an interesting finite dimensional moduli space remains. Finally, I would like to discuss the connection between solitons in noncommutative field theory and D-branes. Actually, these solitons are the D-branes of string theory manifested in a field theory. Therefore, their study allows for probing stringy features in the more controlled context of field theory. An example of a stringy application of our discussion of solitons in scalar noncommutative field theory is in the context of tachyon condensation. It is well-known that bosonic string theory is unstable because of the presence of a tachyonic scalar field $`T`$. The effective action for the tachyon field can be obtained by integrating out massive string fields and is expected to take the form $$S=\frac{C}{g_s}d^{26}x\sqrt{g}\left(\frac{1}{2}f(T)g^{\mu \nu }_\mu T_\nu TV(T)+\mathrm{}\right)$$ (2.48) where higher derivative terms and terms involving massless modes have been neglected. $`V(T)`$ is a general potential with an unstable extremum at $`T=T_0`$ and a minimum chosen to be $`V(0)=0`$. As we have seen in section 1.2.1, turning on a B-field is equivalent to replacing the closed string metric $`g^{\mu \nu }`$ with the effective open string metric $`G^{\mu \nu }`$ and ordinary products with Moyal products in (2.48). In the zero-slope limit derivative terms can be neglected. The solitons of the theory obtained in this limit are the noncommutative solitons we studied before. For instance the gaussian solution $`T=T_0\varphi _0(r)`$ localized in two of the noncommutative directions is a candidate for the $`D23`$-brane. These solitons display the same instability of the corresponding D-branes, since they correspond to an extremum of $`V(T)`$ that is a maximum. From this brief discussion it should be clear that noncommutative field theories exhibit stringy features, such as D-branes, in the simpler context of a field theory. Therefore, their study can be helpful in the understanding of many string theory issues and tachyon condensation is just one example among these. #### 2.2.2 Noncommutative deformation of integrable field theories As we have seen in chapter 1, bosonic noncommutative field theories display a variety of interesting properties but also problematic features, when time is involved in noncommutativity. In this context, an interesting question is how noncommutativity could affect the dynamics of exactly solvable field theories, as for instance two–dimensional integrable theories. As we have seen in section 2.1.1, a common feature of these systems is that the existence of an infinite chain of local conserved currents is guaranteed by the fact that the equations of motion can be written as zero curvature conditions for a suitable set of covariant derivatives . In some cases, as for example the ordinary sine–Gordon or sigma models, an action is also known which generates the integrable equations according to an action principle. Constructing a consistent noncommutative generalization of a two-dimensional theory is a particularly challenging problem though, since, when working with a Minkowski signature, time must be necessarily involved in noncommutativity. Noncommutative versions of ordinary models are intuitively defined as models which reduce to the ordinary ones when the noncommutation parameter $`\theta `$ is removed. As we discussed in detail for the specific example of the free massless scalar field theory in section 1.1.3, in general noncommutative generalizations are not unique as one can construct different noncommutative equations of motion which collapse to the same expression when $`\theta `$ goes to zero. For two dimensional integrable systems, a general criterion to restrict the number of possible noncommutative versions is to require classical integrability to survive in noncommutative geometry. This suggests that any noncommutative generalization should be performed at the level of equations of motion by promoting the standard zero curvature techniques. This program has been worked out for a number of known integrable equations in Refs. . In chapter 1 we have discussed how noncommutative theories naturally arise in the context of string theory. In particular, in section 1.2.2, we have shown how the open $`N=2`$ string in the presence of a constant NS-NS background and a stack of $`n`$ D3-branes can be described, in the zero-slope limit, by $`U(n)`$ noncommutative selfdual Yang-Mills theory. Tree-level S-matrix computations show that the vanishing of amplitudes beyond three point, which is characteristic of ordinary selfdual Yang-Mills, is preserved in the noncommutative theory, suggesting that noncommutative selfdual Yang-Mills, as its ordinary counterpart, is endowed with classical integrability. In section 2.1.3, we have seen that many ordinary integrable models in two and three dimensions can be obtained through a dimensional reduction procedure from selfdual Yang-Mills. From all this it is clear that dimensional reduction from noncommutative selfdual Yang-Mills could be another useful technique to generate possibly integrable noncommutative systems in 1+1 and 2+1 dimensions. It is well known that in integrable commutative field theories there is no particle production and the S-matrix factorizes. A priori the same relation between the existence of infinite conserved charges and factorization properties of scattering processes might be lost in the noncommutative case. Nonlocality in time, responsible for acausal behavior of scattering processes and non unitarity, may interfere in a way to spoil these nice scattering properties. On the other hand, one may also hope that classical integrability would alleviate these pathologies arising when time-space noncommutativity is present. In any case it would be nice to construct a noncommutative generalization of a given ordinary integrable theory characterized by a well-defined and factorized S-matrix. Finally, as we have seen in section 2.1.5, two-dimensional integrable field theories admit soliton solutions. In the noncommutative case, as we have seen in section 2.2.5, a new kind of soliton appears that vanishes in the commutative limit. The class of soliton solutions of the noncommutative version of an integrable field theory is expected to display both solitons that reduce to ordinary ones in the commutative limit and new solitons that vanish in the limit. #### 2.2.3 The natural noncommutative generalization of the sine-Gordon model In my papers , in collaboration with M.T. Grisaru, O. Lechtenfeld, L. Mazzanti, S. Penati and A.D. Popov, continuing the program initiated by M.T. Grisaru and S. Penati in , I have addressed the problem of generalizing the sine-Gordon theory to noncommutative space. The main motivation for this work was the evidence that the natural deformation of this theory, described by the action $$S=\frac{1}{\pi \lambda ^2}d^2z\left[\mathrm{\Phi }\overline{}\mathrm{\Phi }2\gamma (\mathrm{cos}_{}\mathrm{\Phi }1)\right]$$ (2.49) with the corresponding equations of motion $$\overline{}\mathrm{\Phi }=\gamma \mathrm{sin}_{}\mathrm{\Phi }$$ (2.50) is affected by some problems both at the classical and the quantum level. At the classical level it does not seem to be integrable since the ordinary currents promoted to noncommutative currents by replacing the products with $``$–products are not conserved . Moreover, we don’t know how to find a systematic procedure to construct conserved currents since the equations of motion cannot be obtained as zero curvature conditions (a discussion about the lack of integrability for this system is also given in Ref. ). Scattering properties of the natural generalization of the sine-Gordon model have been investigated in . It was found that particle production occurs. The tree level $`24`$ amplitude does not vanish. At the quantum level the renormalizability properties of the ordinary model (2.49) defined for $`\lambda ^2<4`$ seem to be destroyed by noncommutativity. The reason is quite simple and can be understood by analyzing the structure of the divergences of the NC model compared to the ordinary ones . In the $`\lambda ^2<4`$ regime the only divergences come from multitadpole diagrams. In the ordinary case the $`n`$–loop diagram gives a contribution $`(\mathrm{log}m^2a^2)^n`$ where $`a`$ and $`m`$ are the UV and IR cut–offs respectively. This result is independent of the number $`k`$ of external fields and of the external momenta. As a consequence the total contribution at this order can be resummed as $`\gamma (\mathrm{log}m^2a^2)^n(\mathrm{cos}\mathrm{\Phi }1)`$ and the divergence is cancelled by renormalizing the coupling $`\gamma `$. This holds at any order $`n`$ and the model is renormalizable. In the noncommutative case the generic vertex from the expansion of $`\mathrm{cos}_{}\mathrm{\Phi }`$ brings nontrivial phase factors which depend on the momenta coming out of the vertex and on the noncommutativity parameter. The final configuration of phase factors associated to a given diagram depends on the order we use to contract the fields in the vertex. Therefore, in the noncommutative case the ordinary $`n`$–loop diagram splits into a planar and a certain number of nonplanar configurations, where the planar one has a trivial phase factor whereas the nonplanar diagrams differ by the configuration of the phases (for a general discussion see Refs. ). The most general noncommutative multitadpole diagram is built up by combining planar parts with nonplanar ones where two or more tadpoles are intertwined among themselves or with external legs. Since the nonplanar subdiagrams are convergent a generic $`n`$–loop diagram contributes to the divergences of the theory only if it contains a nontrivial planar subdiagram. However, different $`n`$–loop diagrams with different configurations of planar and nonplanar parts give divergent contributions whose coefficients depend on the number $`k`$ of external fields and on the external momenta. A resummation of the divergences to produce a cosine potential is not possible anymore and the renormalization of the couplings of the model is not sufficient to make the theory finite at any order. Noncommutativity seems to deform the cosine potential at the quantum level and the theory loses the renormalizability properties of the corresponding commutative model. Thus, the “natural” generalization of sine–Gordon is not satisfactory and one must look for a different noncommutative generalization compatible with integrability and/or renormalizability. #### 2.2.4 A noncommutative version of the sine-Gordon equation with an infinite number of conserved currents In Ref. M.T. Grisaru and S. Penati constructed a classically integrable noncommutative generalization of the sine-Gordon model, by implementing the bicomplex approach described in section 2.1.1 (as in that section, here we use euclidean signature and complex coordinates $`z=\frac{1}{\sqrt{2}}(x^0+ix^1)`$, $`\overline{z}=\frac{1}{\sqrt{2}}(x^0ix^1)`$). The bicomplex $`(,d,\delta )`$ is considered, where in this case $`=_0L`$, $`_0`$ is the space of $`2\times 2`$ matrices with entries in the algebra of smooth functions on noncommutative $`^2`$ and $`L=_{i=0}^2L^i`$ is a two-dimensional graded vector space with the $`L^1`$ basis $`(\tau ,\sigma )`$ satisfying $`\tau ^2=\sigma ^2=\{\tau ,\sigma \}=0`$. The differential maps are given by $$\delta f=\overline{}f\tau Rf\sigma ;df=Sf\tau +f\sigma $$ (2.51) with commuting constant matrices $`R`$ and $`S`$. The conditions $`d^2=\delta ^2=\{d,\delta \}=0`$ are trivially satisfied. As we have seen in section 2.1.1, to obtain nontrivial second order differential equations the bicomplex must be gauged. In the derivation of the noncommutative sine-Gordon given in the $`d`$ operator is dressed as follows $$D_dfG^1d(Gf)$$ (2.52) with $`G`$ generic invertible matrix in $`_0`$ to be determined in a way to obtain a generalization of the sine-Gordon. While the condition $`D_d^2=0`$ is trivially satisfied, $`\{\delta ,D_d\}=0`$ gives rise to nontrivial second order differential equations $$\overline{}\left(G^1G\right)=[R,G^1G]_{}$$ (2.53) With the choice $`R=S=2\sqrt{\gamma }\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)`$ (2.54) $`G=e_{}^{\frac{\mathrm{i}}{2}\sigma _2\mathrm{\Phi }}=\left(\begin{array}{cc}\mathrm{cos}_{}\frac{\mathrm{\Phi }}{2}& \mathrm{sin}_{}\frac{\mathrm{\Phi }}{2}\\ \mathrm{sin}_{}\frac{\mathrm{\Phi }}{2}& \mathrm{cos}_{}\frac{\mathrm{\Phi }}{2}\end{array}\right)`$ (2.55) equation (2.52) is a matrix equation in $`U(2)`$, corresponding to the system of two coupled equations of motion $`2i\overline{}b\overline{}\left(e_{}^{\frac{i}{2}\mathrm{\Phi }}e_{}^{\frac{i}{2}\mathrm{\Phi }}e_{}^{\frac{i}{2}\mathrm{\Phi }}e_{}^{\frac{i}{2}\mathrm{\Phi }}\right)=i\gamma \mathrm{sin}_{}\mathrm{\Phi }`$ $`2\overline{}a\overline{}\left(e_{}^{\frac{i}{2}\mathrm{\Phi }}e_{}^{\frac{i}{2}\mathrm{\Phi }}+e_{}^{\frac{i}{2}\mathrm{\Phi }}e_{}^{\frac{i}{2}\mathrm{\Phi }}\right)=0`$ (2.56) The first equation contains the potential term which is the “natural” generalization of the ordinary sine potential, whereas the other one has the structure of a conservation law and can be seen as imposing an extra condition on the system. In the commutative limit, the first equation reduces to the ordinary sine–Gordon equation, whereas the second one becomes trivial. The equations are in general complex and possess the $`Z_2`$ symmetry of the ordinary sine–Gordon ( invariance under $`\mathrm{\Phi }\mathrm{\Phi }`$). The reason why integrability seems to require two equations of motions can be traced back to the general structure of unitary groups in noncommutative geometry. In the bicomplex approach the ordinary equations are obtained as zero curvature conditions for covariant derivatives defined in terms of $`SU(2)`$ gauge connections. If the same procedure is to be implemented in the noncommutative case, the group $`SU(2)`$, which is known to be not closed in noncommutative geometry, has to be extended to a noncommutative $`U(2)`$ group and a noncommutative $`U(1)`$ factor enters necessarily into the game. The appearance of the second equation in (2.56) for this noncommutative integrable version of sine-Gordon is then a consequence of the fact that the fields develop a nontrivial trace part. We note that the pattern of equations found in seems to be quite general and unavoidable if integrability is of concern. In fact, the same has been found in Ref. where a different but equivalent set of equations was proposed. In classical integrability of the system described by the set of equations (2.56) was proven by extracting from the bicomplex chain a set of conserved currents that are local, in the sense that they are functions of the field $`\mathrm{\Phi }`$ and its derivatives, but not of integrals of $`\mathrm{\Phi }`$. Since Moyal product has an infinite expansion in derivatives of fields, this introduces a kind of nonlocality in the theory that is intrinsic and unavoidable when working in a noncommutative space. In the expansion in the noncommutativity parameter $`\theta `$ has been studied for the first currents to check their relation with ordinary currents and to explicitly verify their conservation up to second order in the deformation parameter. #### 2.2.5 Solitons The presence of two equations of motion is in principle very restrictive and one may wonder whether the class of solutions is empty. To show that this is not the case, in Ref. solitonic solutions were constructed perturbatively which reduce to the ordinary solitons when we take the commutative limit. Since a classical action was not found in , these solitonic solutions found are said to be localized in the sense that at order zero in the deformation parameter they reduce to the well-known euclidean solitons of the sine-Gordon theory. Since the solution at order zero determines the solution to all orders in the deformation parameter, in these solitons are called localized at all orders. More generally, we observe that the second equation in (2.56) is automatically satisfied by any chiral or antichiral function. Therefore, we expect the class of solitonic solutions to be at least as large as the ordinary one. In the general case, instead, we expect the class of dynamical solutions to be smaller than the ordinary one because of the presence of the nontrivial constraint. However, since the constraint equation is one order higher with respect to the dynamical equation, order by order in the $`\theta `$-expansion a solution always exists. This means that a Seiberg–Witten map between the NC and the ordinary model does not exist as a mapping between physical configurations, but it might be constructed as a mapping between equations of motion or conserved currents. The kind of noncommutative solitons discussed in section 2.2.1 has not been studied in . These solutions in principle should be present in this model. However, the model described in was shown to display bad scattering properties , as I will show in detail in section 2.2.9. For this reason it had to be discarded and replaced with a new model described in section 2.3 . Both kinds of soliton solutions were studied in detail for this model (see section 2.3.5). #### 2.2.6 Reduction from noncommutative selfdual Yang-Mills The material presented in this section and the following ones, until the end of section 2.2, is mostly taken from the paper , written in collaboration with M.T. Grisaru, L. Mazzanti and S. Penati. The (anti-)selfdual Yang–Mills equation is well-known to describe a completely integrable classical system in four dimensions . In the ordinary case the equations of motion for many two dimensional integrable systems, including sine–Gordon, can be obtained through dimensional reduction of the (anti)selfdual Yang-Mills equations . We have seen in section 2.1.3 that a convenient description of the (anti)selfdual Yang-Mills system is the so called $`J`$-formulation, given in terms of a $`SL(N,)`$ matrix-valued $`J`$ field satisfying $$_{\overline{y}}\left(J^1_yJ\right)+_{\overline{z}}\left(J^1_zJ\right)=0$$ (2.57) where $`y`$, $`\overline{y}`$, $`z`$, $`\overline{z}`$ are complex variables treated as formally independent. In the ordinary case, the sine-Gordon equation can be obtained from (2.57) by taking $`J`$ in $`SL(2,)`$ to be $$J=J(u,z,\overline{z})=e^{\frac{z}{2}\sigma _i}e^{\frac{i}{2}\mathrm{\Phi }\sigma _j}e^{\frac{\overline{z}}{2}\sigma _i}$$ (2.58) where $`\mathrm{\Phi }=\mathrm{\Phi }(y,\overline{y})`$ depends on $`y`$ and $`\overline{y}`$ only and $`\sigma _i`$ are the Pauli matrices. A noncommutative version of the (anti-)selfdual Yang–Mills system can be naturally obtained by promoting the variables $`y`$, $`\overline{y}`$, $`z`$ and $`\overline{z}`$ to be noncommutative thus extending the ordinary products in (2.57) to $``$–products. In this case the $`J`$ field lives in $`GL(N,)`$. As we outlined in section 1.2.3, it has been shown that noncommutative selfdual Yang-Mills naturally emerges from open $`N=2`$ strings in a B-field background. Moreover, in examples of reductions to two-dimensional noncommutative systems were given. It was also argued that the noncommutative deformation should preserve the integrability of the systems . We now show that our noncommutative version of the sine-Gordon equations can be derived through dimensional reduction from the noncommutative selfdual Yang-Mills equations. For this purpose we consider the noncommutative version of equations (2.57) and choose $`J_{}`$ in $`GL(2,)`$ as $$J_{}=J_{}(u,z,\overline{z})=e_{}^{\frac{z}{2}\sigma _i}e_{}^{\frac{i}{2}\mathrm{\Phi }\sigma _j}e_{}^{\frac{\overline{z}}{2}\sigma _i}$$ (2.59) This leads to the matrix equation $$_{\overline{y}}aI+i\left(_{\overline{y}}b+\frac{1}{2}\mathrm{sin}_{}\mathrm{\Phi }\right)\sigma _j=0$$ (2.60) where $`a`$ and $`b`$ have been defined in (2.56). Now, taking the trace we obtain $`_{\overline{y}}a=0`$ which is the constraint equation in (2.56). As a consequence, the term proportional to $`\sigma _j`$ gives rise to the dynamical equation in (2.56) for the particular choice $`\gamma =1`$. Therefore we have shown that the equations of motion of the noncommutative version of sine–Gordon proposed in can be obtained from a suitable reduction of the noncommutative selfdual Yang-Mills system as in the ordinary case. From this derivation the origin of the constraint appears even more clearly: it arises from setting to zero the trace part which the matrices in $`GL(2,)`$ naturally develop under $``$–multiplication. Solving (2.60) for the particular choice $`\sigma _j=\sigma _3`$ we obtain the alternative set of equations $`\overline{}\left(e_{}^{\frac{i}{2}\mathrm{\Phi }}e_{}^{\frac{i}{2}\mathrm{\Phi }}\right)={\displaystyle \frac{i}{2}}\gamma \mathrm{sin}_{}\mathrm{\Phi }`$ (2.61) $`\overline{}\left(e_{}^{\frac{i}{2}\mathrm{\Phi }}e_{}^{\frac{i}{2}\mathrm{\Phi }}\right)={\displaystyle \frac{i}{2}}\gamma \mathrm{sin}_{}\mathrm{\Phi }`$ (2.62) Order by order in the $`\theta `$-expansion the set of equations (2.56) and (2.62) are equivalent. Therefore, the set (2.62) is equally suitable for the description of an integrable noncommutative generalization of sine-Gordon. Since our noncommutative generalization of sine–Gordon is integrable, the present result gives support to the arguments in favor of the integrability of noncommutative selfdual Yang-Mills system. #### 2.2.7 The action We are now interested in the possibility of determining an action for the scalar field $`\mathrm{\Phi }`$ satisfying the system of eqs. (2.56). We are primarily motivated by the possibility to move on to a quantum description of the system. In general, it is not easy to find an action for the dynamical equation (the first eq. in (2.56)) since $`\mathrm{\Phi }`$ is constrained by the second one. One possibility could be to implement the constraint by the use of a Lagrange multiplier. Another quite natural possibility is to try to obtain the action by a dimensional reduction procedure from $`(4,0)`$ selfdual Yang-Mills action in Yang formulation. Unfortunately, this does not work, since WZW-like terms disappear from the reduced action because of cyclicity of Moyal product in an integral. As a result one obtains a reduced action generating nonchiral equations, different from the chiral ones in eqs. (2.56) and (2.62). We consider instead the equivalent set of equations (2.62). We rewrite them in the form $`\overline{}(g^1g)={\displaystyle \frac{1}{4}}\gamma \left(g^2g^2\right)`$ (2.63) $`\overline{}(gg^1)={\displaystyle \frac{1}{4}}\gamma \left(g^2g^2\right)`$ (2.64) where we have defined $`ge_{}^{\frac{i}{2}\mathrm{\Phi }}`$. Since $`\mathrm{\Phi }`$ is in general complex $`g`$ can be seen as an element of a noncommutative complexified $`U(1)`$. The gauge group valued function $`\overline{g}(g^{})^1=e_{}^{\frac{i}{2}\mathrm{\Phi }^{}}`$ is subject to the equations $`\overline{}(\overline{g}\overline{g}^1)={\displaystyle \frac{1}{4}}\gamma \left(\overline{g}^2\overline{g}^2\right)`$ (2.65) $`\overline{}(\overline{g}^1\overline{g})={\displaystyle \frac{1}{4}}\gamma \left(\overline{g}^2\overline{g}^2\right)`$ (2.66) obtained by taking the h.c. of (2.64). In order to determine the action it is convenient to concentrate on the first equation in (2.64) and the second one in (2.66) as the two independent complex equations of motion which describe the dynamics of our system. We first note that the left-hand sides of equations (2.64) and (2.66) have the chiral structure which is well known to correspond to a noncommutative version of the WZNW action (see section 1.1.3). Therefore we are led to consider the action $$S[g,\overline{g}]=S[g]+S[\overline{g}]$$ (2.67) where, introducing the homotopy path $`\widehat{g}(t)`$ such that $`\widehat{g}(0)=1`$, $`\widehat{g}(1)=g`$ ($`t`$ is a commuting parameter) we have defined $`S[g]`$ $`={\displaystyle }d^2z[g\overline{}g^1+{\displaystyle _0^1}dt\widehat{g}^1_t\widehat{g}[\widehat{g}^1\widehat{g},\widehat{g}^1\overline{}\widehat{g}]_{}`$ (2.68) $``$ $`{\displaystyle \frac{\gamma }{4}}(g^2+g^22)]`$ (2.69) and similarly for $`S[\overline{g}]`$. The first part of the action can be recognized as the noncommutative generalization of a complexified $`U(1)`$ WZNW action . To prove that this generates the correct equations, we should take the variation with respect to the $`\mathrm{\Phi }`$ field ($`g=e_{}^{\frac{i}{2}\mathrm{\Phi }}`$) and deal with complications which follow from the fact that in the noncommutative case the variation of an exponential is not proportional to the exponential itself. However, since the variation $`\delta \mathrm{\Phi }`$ is arbitrary, we can forget about its $`\theta `$ dependence and write $`\frac{i}{2}\delta \mathrm{\Phi }=g^1\delta g`$, trading the variation with respect to $`\mathrm{\Phi }`$ with the variation with respect to $`g`$. Analogously, the variation with respect to $`\mathrm{\Phi }^{}`$ can be traded with the variation with respect to $`\overline{g}`$. It is then a simple calculation to show that $$\delta S[g]=d^2z2g^1\delta g\left[\overline{}\left(g^1g\right)\frac{i}{2}\gamma \mathrm{sin}_{}\mathrm{\Phi }\right]$$ (2.70) from which we obtain the first equation in (2.64). Treating $`\overline{g}`$ as an independent variable an analogous derivation gives the second equation in (2.66) from $`S[\overline{g}]`$. We note that, when $`\mathrm{\Phi }`$ is real, $`g=\overline{g}`$ and the action (2.67) reduces to $`S_{\mathrm{real}}[g]=2S_{WZW}[g]\gamma (\mathrm{cos}_{}\mathrm{\Phi }1)`$. In general, since the two equations (2.56) are complex it would be inconsistent to restrict ourselves to real solutions. However, it is a matter of fact that the equations of motion become real when the field is real. Perturbatively in $`\theta `$ this can be proved order by order by direct inspection of the equations in . In particular, at a given order one can show that the imaginary part of the equations vanishes when the constraint and the equations of motion at lower orders are satisfied. #### 2.2.8 The relation to the noncommutative Thirring model In the ordinary case the equivalence between the Thirring and sine–Gordon models can be proven at the level of functional integrals by implementing the bosonization prescription on the fermions. The same procedure has been worked out in noncommutative geometry . Starting from the noncommutative version of Thirring described by $$S_T=d^2x\left[\overline{\psi }i\gamma ^\mu _\mu \psi +m\overline{\psi }\psi \frac{\lambda }{2}(\overline{\psi }\gamma ^\mu \psi )(\overline{\psi }\gamma _\mu \psi )\right]$$ (2.71) the bosonization prescription gives rise to the action for the bosonized noncommutative massive Thirring model which turns out to be a noncommutative WZNW action supplemented by a cosine potential term for the noncommutative U(1) group valued field which enters the bosonization of the fermionic currents. In particular, in the most recent paper in Ref. it has been shown that working in Euclidean space the massless Thirring action corresponds to the sum of two WZNW actions once a suitable choice for the regularization parameter is made. Moreover, in Ref. it was proven that the bosonization of the mass term in (2.71) gives rise to a cosine potential for the scalar field with coupling constant proportional to $`m`$. The main observation is that our action (2.67) is the sum of two noncommutative WZNW actions plus cosine potential terms for the pair of $`U(1)_{}`$ group valued fields $`g`$ and $`\overline{g}`$, considered as independent. Therefore, our action can be interpreted as coming from the bosonization of the massive noncommutative Thirring model, in agreement with the results in . We have shown that even in the noncommutative case the sine–Gordon field can be interpreted as the scalar field which enters the bosonization of the Thirring model, so proving that the equivalence between the Thirring and sine–Gordon models can be maintained in noncommutative generalizations of these models. Moreover, the classical integrability of our noncommutative version of sine–Gordon proven in should automatically guarantee the integrability of the noncommutative Thirring model. In the particular case of zero coupling ($`\gamma =0`$), the equations (2.67) and (2.62) correspond to the action and the equations of motion for a noncommutative $`U(1)`$ WZNW model , respectively. Again, we can use the results of to prove the classical integrability of the noncommutative $`U(1)`$ WZNW model and construct explicitly its conserved currents. #### 2.2.9 (Bad) properties of the S-matrix In section 2.1.4 we showed that in integrable commutative field theories there is no particle production and the S-matrix factorizes. In the noncommutative case properties of the S-matrix have been investigated for two specific models: The $`\lambda \mathrm{\Phi }^4`$ theory in two dimensions and the nonintegrable “natural” NC generalization the the sine–Gordon model . In the first reference a very pathological acasual behavior was observed due to the space and time noncommutativity (see section 1.1.2). For an incoming wave packet the scattering produces an advanced wave which arrives at the origin before the incoming wave. In the second model investigated it was found that particle production occurs. The tree level $`24`$ amplitude does not vanish. It might be hoped that classical integrability would alleviate these pathologies. In the NC integrable sine-Gordon case, since we have an action, it is possible to investigate these issues. As described below we have computed the scattering amplitude for the $`22`$ process and found that the acausality of Ref. is not cured by integrability. We have also computed the production amplitudes for the processes $`23`$ and $`24`$ and found that they don’t vanish. We started from our action (2.69) rewritten in terms of Minkowski space coordinates $`x^0,x^1`$ and real fields ($`g=e_{}^{\frac{i}{2}\mathrm{\Phi }}`$, $`\widehat{g}(t)=e_{}^{\frac{i}{2}t\mathrm{\Phi }}`$ with $`\mathrm{\Phi }`$ real) $`S[g]={\displaystyle \frac{1}{2}}{\displaystyle d^2xg^1^\mu gg^1_\mu g}+{\displaystyle \frac{\gamma }{4}}{\displaystyle d^2x(g^2+g^22)}`$ (2.72) $`{\displaystyle \frac{1}{3}}{\displaystyle d^3xϵ^{\mu \nu \rho }\widehat{g}^1_\mu \widehat{g}\widehat{g}^1_\nu \widehat{g}\widehat{g}^1_\rho \widehat{g}}`$ (2.73) where $`fg=fe^{\frac{i}{2}\theta ϵ^{\mu \nu }\stackrel{}{}_\mu \stackrel{}{}_\nu }g`$, and we derived the following Feynman’s rules * The propagator $$G(q)=\frac{4i}{q^22\gamma }$$ (2.74) * The vertices $`v_3(k_1,\mathrm{},k_3)={\displaystyle \frac{2}{2^33!}}ϵ^{\mu \nu }k_{1\mu }k_{2\nu }F(k_1,\mathrm{},k_3)`$ $`v_4(k_1,\mathrm{},k_4)=i\left({\displaystyle \frac{1}{2^44!}}\left(k_1^2+3k_1k_3\right)+{\displaystyle \frac{\gamma }{24!}}\right)F(k_1,\mathrm{},k_4)`$ $`v_5(k_1,\mathrm{},k_5)={\displaystyle \frac{2ϵ^{\mu \nu }}{2^55!}}\left(k_{1\mu }k_{2\nu }k_{1\mu }k_{3\nu }+2k_{1\mu }k_{4\nu }\right)F(k_1,\mathrm{},k_5)`$ $`v_6(k_1,\mathrm{},k_6)=i[{\displaystyle \frac{1}{2^66!}}(k_1^2+5k_1k_35k_1k_4+5k_1k_5){\displaystyle \frac{\gamma }{26!}}]`$ (2.75) $`F(k_1,\mathrm{},k_6)`$ (2.76) where $$F(k_1,\mathrm{},k_n)=\mathrm{exp}\left(\frac{i}{2}\underset{i<j}{}k_i\times k_j\right)$$ (2.77) is the phase factor coming from the $``$-products in the action (we have indicated $`a\times b=\theta ϵ^{\mu \nu }a_\mu b_\nu `$), $`k_i`$ are all incoming momenta and we used momentum conservation. At tree level the $`22`$ process is described by the diagrams with the topologies in Fig. 1. Figure 1: Tree level $`22`$ amplitude Including contributions from the various channels and using the three point and four point vertices of eqs. (2.76) we obtained for the scattering amplitude the expression $$\frac{i}{2}E^2p^2\left(\frac{1}{2E^2\gamma }\frac{1}{2p^2+\gamma }\right)\mathrm{sin}^2(pE\theta )+i\frac{\gamma }{2}\mathrm{cos}^2(pE\theta )$$ (2.78) where $`p`$ is the center of mass momentum and $`E=\sqrt{p^2+2\gamma }`$. For comparison with Ref. this should be multiplied by an incoming wave packet $$\varphi _{in}(p)\left(e^{\frac{(pp_0)^2}{\lambda }}+e^{\frac{(p+p_0)^2}{\lambda }}\right)$$ (2.79) and Fourier transformed with $`e^{ipx}`$. We have not attempted to carry out the Fourier transform integration. However, we note that for $`p_0`$ very large $`E`$ and $`p`$ are concentrated around large values and the scattering amplitude assumes the form $$i\frac{\gamma }{4}\mathrm{sin}^2(pE\theta )+i\frac{\gamma }{2}\mathrm{cos}^2(pE\theta )$$ (2.80) which is equivalent to the result in Ref. , leading to the same acausal pathology <sup>1</sup><sup>1</sup>1It is somewhat tantalizing that a change in the relative coefficient between the two terms would lead to a removal of the trigonometric factors which are responsible for the acasual behavior.. We describe now the computation of the production amplitudes $`23`$ and $`24`$. At tree level the contributions are drawn in Figures 2 and 3, respectively. Figure 2: Tree level $`23`$ amplitude Figure 3: Tree level $`24`$ amplitude For any topology the different possible channels must be taken into account. This, as well as the complicated expressions for the vertices, has led us to use an algebraic manipulation program computer. We used Mathematica<sup>©</sup> to symmetrize completely the vertices (2.76). This allows to take automatically into account the different diagrams obtained by exchanging momenta entering a given vertex. The contribution from each diagram was obtained as a product of the combinatorial factor, the relevant vertices and propagators. Due to the length of the program it was impossible to handle the calculation in a completely analytic way. Instead, the program was run with assigned values of the momenta and arbitrary $`\theta `$ and $`\gamma `$. For both the $`23`$ and $`24`$ processes the result is nonvanishing. As a check of our calculation we mention that the production amplitudes vanish when we set $`\theta =0`$, for any value of the coupling and the momenta. #### 2.2.10 Conclusions In , in collaboration with M.T. Grisaru, L. Mazzanti and S. Penati I have investigated some properties of the integrable noncommutative sine–Gordon system proposed in . We succeeded in constructing an action which turned out to be a WZNW action for a noncommutative, complexified $`U(1)`$ augmented by a cosine potential. We have shown that even in the noncomutative case there is a duality relation between our integrable noncommutative sine–Gordon model and the noncommutative Thirring. Noncommutative WZNW models have been shown to be one–loop renormalizable . This suggests that the noncommutative sine–Gordon model proposed in is not only integrable but it might lead to a well-defined quantized model, giving support to the existence of a possible relation between integrability and renormalizability. Armed with our action we investigated some properties of the S–matrix for elementary excitations. However, in contradistinction to the commutative case, the S–matrix turned out to be acasual and nonfactorizable <sup>2</sup><sup>2</sup>2Other problems of the S-matrix have been discussed in .. The reason for the acasual behavior has been discussed in where it was pointed out that noncommutativity induces a backward-in-time effect because of the presence of certain phase factors (see section 1.1.2). It appears that in our case this effect is still present in spite of integrability. It is not clear why the presence of an infinite number of local conserved currents (local in the sense that they are not expressed as integrals of certain densities) does not guarantee factorization and absence of production in the S-matrix as it does in the commutative case. The standard proofs of factorization use, among other assumptions, the mutual commutativity of the charges - a property we have not been able to check so far because of the complicated nature of the currents. But even if the charges were to commute the possibility of defining them as powers of momenta, as required in the proofs, could be spoiled by acausal effects which prevent a clear distinction between incoming and outgoing particles. In a series of papers a different approach to quantum noncommutative field theories has been proposed when the time variable is not commuting. In those papers it has been argued that the problems associated to time-space noncommutativity are due to the fact that the time-ordering procedure does not commute with the star multiplication. Starting from the usual definition of the S-matrix in terms of the time-ordered exponential of the interaction term in the action and applying Wick theorem, one cannot combine the contraction functions of positive and negative frequency to obtain the causal Feynman propagator. Therefore, it has been suggested that, instead of the Feyman approach , one should use the time ordered perturbation theory extended to the noncommutative case. Moreover, it has been shown that in this framework unitarity is preserved as long as the interaction lagrangian is explicitly hermitian. It would be interesting to redo our calculations in that approach to see whether a well-defined factorized S-matrix for our model can be constructed. In this context it would be also interesting to investigate the scattering of solitons present in our model . In the next section a novel noncommutative sine-Gordon system, obtained by dimensional reduction from the $`2+1`$ model introduced in , will be constructed and studied. We will see that it exhibits nice scattering properties, consistent with the usual relation between integrability and factorization of the S-matrix. ### 2.3 The noncommutative integrable sine-Gordon model #### 2.3.1 Introduction In this section I will present the results I obtained, in collaboration with O. Lechtenfeld, L. Mazzanti, S. Penati and A.D. Popov, in . The main goal of that work was to find a noncommutative generalization of the sine-Gordon system which, as a hallmark of integrability, possesses a well-defined causal and factorized S-matrix. Furthermore, its equations of motion should admit noncommutative multi-soliton solutions which represent deformations of the well known sine-Gordon solitons. In sections 2.2.4 and 2.2.5 I have discussed the results obtained in , where a model was proposed which describes the dynamics of a complex scalar field by a couple of equations of motion. These equations were obtained as flatness conditions for a $`U(2)`$ bidifferential calculus and automatically guarantee the existence of an infinite number of local conserved currents. The same equations were also generated in via a particular dimensional reduction of the noncommutative $`U(2)`$ selfdual Yang-Mills equations in euclidean space. However, this reduction did not work at the level of the action, which turned out to be the sum of two WZW models augmented by a cosine potential. Evaluating tree-level scattering amplitudes it was discovered, furthermore, that this model suffers from acausal behavior and a non-factorized S-matrix, meaning that particle production occurs. At this point it is important to note that the noncommutative deformation of an integrable equation is a priori not unique, because one may always add terms which vanish in the commutative limit, as we have seen in section 1.1.3. For the case at hand, for example, different inequivalent ansätze for the $`U(2)`$ matrices entering the bicomplex construction are possible as long as they all reproduce the ordinary sine-Gordon equation in the commutative limit. It is therefore conceivable that among these possibilities there exists an ansatz (different from the one in ) which guarantees the classical integrability of the corresponding noncommutative model. What is already certain is the necessity to introduce two real scalar fields instead of one, since in the noncommutative realm the $`U(1)`$ subgroup of $`U(2)`$ fails to decouple. What has been missing is a guiding principle towards the “correct” field parametrization. Since the sine-Gordon model can be obtained by dimensional reduction from 2+2 dimensional selfdual Yang-Mills theory via a 2+1 dimensional integrable sigma model , and because the latter’s noncommutative extension was shown to be integrable in , it seems a good idea to contruct an integrable generalization of the sine-Gordon equation by starting from the linear system of this integrable sigma model endowed with a time-space noncommutativity. This is the key strategy of this paper. The reduction is performed on the equations of motion first, but it also works at the level of the action, so giving directly the 1+1 dimensional action we are looking for. This success is an indication that the new field parametrization proposed in is the proper one. To be more precise, in we proposed three different parametrizations, by pairs of fields $`(\varphi _+,\varphi _{})`$, $`(\rho ,\phi )`$ and $`(h_1,h_2)`$, all related by nonlocal field redefinitions but all deriving from the compatibility conditions of the underlying linear system . The first two appear in Yang formulation while the third one arises in Leznov formulation . For either field pair in Yang formulation, the nontrivial compatibility condition reduces to a pair of “noncommutative sine-Gordon equations” which in the commutative limit degenerates to the standard sine-Gordon equation for $`\frac{1}{2}(\varphi _++\varphi _{})`$ or $`\phi `$, respectively, while $`\frac{1}{2}(\varphi _+\varphi _{})`$ or $`\rho `$ decouple as free bosons. The alternative Leznov formulation has the advantage of producing two polynomial (actually, quadratic) equations of motion for $`(h_1,h_2)`$ but retains their coupling even in the commutative limit. With the linear system comes a well-developed technology for generating solitonic solutions to the equations of motion. In the dressing method was employed to explicitly outline the construction of noncommutative sine-Gordon multi-solitons, directly in 1+1 dimensions as well as by reducing plane-wave solutions of the 2+1 dimensional integrable sigma model . The one-soliton sector was completely analyzed and it was found that the standard soliton solution are recovered as undeformed. Noncommutativity becomes palpable only at the multi-soliton level. It was shown in that the tree-level $`n`$-point amplitudes of noncommutative 2+2 dimensional SDYM vanish for $`n>3`$, consistent with the vanishing theorems for the $`N=2`$ string. Therefore, we were expecting nice properties of the S-matrix to be inherited by this noncommutative sine-Gordon theory. Indeed, a direct evaluation of tree-level amplitudes revealed that, in the Yang as well as the Leznov formulation, the S-matrix is causal and no particle production occurs. #### 2.3.2 A noncommutative integrable sigma model in $`2+1`$ dimensions As has been known for some time, nonlinear sigma models in $`2+1`$ dimensions may be Lorentz-invariant or integrable but not both . Since the integrable variant, introduced in section 2.1.3, serves as our starting point for the derivation of the sine-Gordon model and its soliton solutions, we shall present its noncommutative extension in some detail in the present section. ##### Conventions in noncommutative $`^{2,1}`$ In $`^{2,1}`$ we shall use (real) coordinates $`(x^a)=(t,x,y)`$ in which the Minkowskian metric reads $`(\eta _{ab})=\text{diag}(1,+1,+1)`$. For later use we introduce the light-cone coordinates $$u:=\frac{1}{2}(t+y),v:=\frac{1}{2}(ty),_u=_t+_y,_v=_t_y.$$ (2.81) In view of the future reduction to $`1+1`$ dimensions, we choose the coordinate $`x`$ to remain commutative, so that the only non-vanishing component of the noncommutativity tensor is $$\theta ^{ty}=\theta ^{yt}=:\theta >0.$$ (2.82) ##### Linear system Consider on noncommutative $`^{2,1}`$ the following pair of linear differential equations , $$(\zeta _x_u)\mathrm{\Psi }=A\mathrm{\Psi }\text{and}(\zeta _v_x)\mathrm{\Psi }=B\mathrm{\Psi },$$ (2.83) where a spectral parameter $`\zeta P^1S^2`$ has been introduced. The auxiliary field $`\mathrm{\Psi }`$ takes values in U$`(n)`$ and depends on $`(t,x,y,\zeta )`$ or, equivalently, on $`(x,u,v,\zeta )`$. The $`u(n)`$ matrices $`A`$ and $`B`$, in contrast, do not depend on $`\zeta `$ but only on $`(x,u,v)`$. Given a solution $`\mathrm{\Psi }`$, they can be reconstructed via<sup>3</sup><sup>3</sup>3 Inverses are understood with respect to the star product, i.e. $`\mathrm{\Psi }^1\mathrm{\Psi }=\mathrm{𝟏}`$. $$A=\mathrm{\Psi }(_u\zeta _x)\mathrm{\Psi }^1\text{and}B=\mathrm{\Psi }(_x\zeta _v)\mathrm{\Psi }^1.$$ (2.84) It should be noted that the equations (2.83) are not of first order but actually of infinite order in derivatives, due to the star products involved. In addition, the matrix $`\mathrm{\Psi }`$ is subject to the following reality condition : $$\mathrm{𝟏}=\mathrm{\Psi }(t,x,y,\zeta )[\mathrm{\Psi }(t,x,y,\overline{\zeta })]^{},$$ (2.85) where ‘$``$’ is hermitian conjugation. The compatibility conditions for the linear system (2.83) read $`_xB_vA=0,`$ (2.86) $`_xA_uBAB+BA=0.`$ (2.87) By detailing the behavior of $`\mathrm{\Psi }`$ at small $`\zeta `$ and at large $`\zeta `$ we shall now “solve” these equations in two different ways, each one leading to a single equation of motion for a particular field theory. ##### Yang-type solution We require that $`\mathrm{\Psi }`$ is regular at $`\zeta =0`$ , $$\mathrm{\Psi }(t,x,y,\zeta 0)=\mathrm{\Phi }^1(t,x,y)+O(\zeta ),$$ (2.88) which defines a U$`(n)`$-valued field $`\mathrm{\Phi }(t,x,y)`$, i.e. $`\mathrm{\Phi }^{}=\mathrm{\Phi }^1`$. Therewith, $`A`$ and $`B`$ are quickly reconstructed via $$A=\mathrm{\Psi }_u\mathrm{\Psi }^1\|_{\zeta =0}=\mathrm{\Phi }^1_u\mathrm{\Phi }\text{and}B=\mathrm{\Psi }_x\mathrm{\Psi }^1\|_{\zeta =0}=\mathrm{\Phi }^1_x\mathrm{\Phi }$$ (2.89) It is easy to see that compatibility equation (2.87) is then automatic while the remaining equation (2.86) turns into $$_x(\mathrm{\Phi }^1_x\mathrm{\Phi })_v(\mathrm{\Phi }^1_u\mathrm{\Phi })=0.$$ (2.90) This Yang-type equation can be rewritten as $$(\eta ^{ab}+v_cϵ^{cab})_a(\mathrm{\Phi }^1_b\mathrm{\Phi })=0,$$ (2.91) where $`ϵ^{abc}`$ is the alternating tensor with $`ϵ^{012}=1`$ and $`(v_c)=(0,1,0)`$ is a fixed spacelike vector. Clearly, this equation is not Lorentz-invariant but (deriving from a Lax pair) it is integrable. One can recognize (2.91) as the field equation for (a noncommutative generalization of) a WZW-like modified U$`(n)`$ sigma model with the action<sup>4</sup><sup>4</sup>4 which is obtainable by dimensional reduction from the Nair-Schiff action for SDYM in 2+2 dimensions $`S_\text{Y}=\frac{1}{2}{\displaystyle dtdxdy\eta ^{ab}\mathrm{tr}\left(_a\mathrm{\Phi }^1_b\mathrm{\Phi }\right)}`$ (2.92) $`\frac{1}{3}{\displaystyle dtdxdy_0^1d\lambda \stackrel{~}{v}_\rho ϵ^{\rho \mu \nu \sigma }\mathrm{tr}\left(\stackrel{~}{\mathrm{\Phi }}^1_\mu \stackrel{~}{\mathrm{\Phi }}\stackrel{~}{\mathrm{\Phi }}^1_\nu \stackrel{~}{\mathrm{\Phi }}\stackrel{~}{\mathrm{\Phi }}^1_\sigma \stackrel{~}{\mathrm{\Phi }}\right)}`$ where Greek indices include the extra coordinate $`\lambda `$, and $`ϵ^{\rho \mu \nu \sigma }`$ denotes the totally antisymmetric tensor in $`^4`$. The field $`\stackrel{~}{\mathrm{\Phi }}(t,x,y,\lambda )`$ is an extension of $`\mathrm{\Phi }(t,x,y)`$, interpolating between $$\stackrel{~}{\mathrm{\Phi }}(t,x,y,0)=\text{const}\text{and}\stackrel{~}{\mathrm{\Phi }}(t,x,y,1)=\mathrm{\Phi }(t,x,y),$$ (2.93) and ‘$`\mathrm{tr}`$’ implies the trace over the U$`(n)`$ group space. Finally, $`(\stackrel{~}{v}_\rho )=(v_c,0)`$ is a constant vector in (extended) space-time. ##### Leznov-type solution Finally, we also impose the asymptotic condition that $`lim_\zeta \mathrm{}\mathrm{\Psi }=\mathrm{\Psi }^0`$ with some constant unitary (normalization) matrix $`\mathrm{\Psi }^0`$. The large $`\zeta `$ behavior $$\mathrm{\Psi }(t,x,y,\zeta \mathrm{})=\left(\mathrm{𝟏}+\zeta ^1\mathrm{{\rm Y}}(t,x,y)+O(\zeta ^2)\right)\mathrm{\Psi }^0$$ (2.94) then defines a $`u(n)`$-valued field $`\mathrm{{\rm Y}}(t,x,y)`$. Again this allows one to reconstruct $`A`$ and $`B`$ through $$A=\underset{\zeta \mathrm{}}{lim}\left(\zeta \mathrm{\Psi }_x\mathrm{\Psi }^1\right)=_x\mathrm{{\rm Y}}\text{and}B=\underset{\zeta \mathrm{}}{lim}\left(\zeta \mathrm{\Psi }_v\mathrm{\Psi }^1\right)=_v\mathrm{{\rm Y}}$$ (2.95) In this parametrization, compatibility equation (2.86) becomes an identity but the second equation (2.87) turns into $$_x^2\mathrm{{\rm Y}}_u_v\mathrm{{\rm Y}}_x\mathrm{{\rm Y}}_v\mathrm{{\rm Y}}+_v\mathrm{{\rm Y}}_x\mathrm{{\rm Y}}=0.$$ (2.96) This Leznov-type equation can also be obtained by extremizing the action $$S_\text{L}=dtdxdy\mathrm{tr}\left\{\frac{1}{2}\eta ^{ab}_a\mathrm{{\rm Y}}_b\mathrm{{\rm Y}}+\frac{1}{3}\mathrm{{\rm Y}}\left(_x\mathrm{{\rm Y}}_v\mathrm{{\rm Y}}_v\mathrm{{\rm Y}}_x\mathrm{{\rm Y}}\right)\right\},$$ (2.97) which is merely cubic. Obviously, the Leznov field $`\mathrm{{\rm Y}}`$ is related to the Yang field $`\mathrm{\Phi }`$ through the non-local field redefinition $$_x\mathrm{{\rm Y}}=\mathrm{\Phi }^1_u\mathrm{\Phi }\text{and}_v\mathrm{{\rm Y}}=\mathrm{\Phi }^1_x\mathrm{\Phi }.$$ (2.98) For each of the two fields $`\mathrm{\Phi }`$ and $`\mathrm{{\rm Y}}`$, one equation from the pair (2.86, 2.87) represents the equation of motion, while the other one is a direct consequence of the parametrization (2.89) or (2.95). #### 2.3.3 Reduction to noncommutative sine-Gordon ##### Algebraic reduction ansatz In section 2.1.3 we have seen that the (commutative) sine-Gordon equation can be obtained from the self-duality equations for SU(2) Yang-Mills upon appropriate reduction from $`2+2`$ to $`1+1`$ dimensions. In this process the integrable sigma model of the previous section appears as an intermediate step in $`2+1`$ dimensions, and so we may take its noncommutative extension as our departure point, after enlarging the group to U(2). In order to avoid cluttering the formulae we suppress the ‘$``$’ notation for noncommutative multiplication from now on: all products are assumed to be star products, and all functions are built on them, i.e. $`e^{f(x)}`$ stands for $`e_{}^{f(x)}`$ and so on. The dimensional reduction proceeds in two steps, firstly, a factorization of the coordinate dependence and, secondly, an algebraic restriction of the form of the U(2) matrices involved. In the language of the linear system (2.83) the adequate ansatz for the auxiliary field $`\mathrm{\Psi }`$ reads $$\mathrm{\Psi }(t,x,y,\zeta )=V(x)\psi (u,v,\zeta )V^{}(x)\text{with}V(x)=e^{\mathrm{i}\alpha x\sigma _1},$$ (2.99) where $`\sigma _1=(\begin{array}{cc}0& 1\\ 1& 0\end{array})`$, $``$ denotes some constant unitary matrix (to be specified later) and $`\alpha `$ is a constant parameter. Under this factorization, the linear system (2.83) simplifies to<sup>5</sup><sup>5</sup>5 The adjoint action means $`\mathrm{ad}\sigma _1(\psi )=[\sigma _1,\psi ]`$. $$(_u\mathrm{i}\alpha \zeta \mathrm{ad}\sigma _1)\psi =a\psi \text{and}(\zeta _v\mathrm{i}\alpha \mathrm{ad}\sigma _1)\psi =b\psi $$ (2.100) with $`a=V^{}AV`$ and $`b=V^{}BV`$. Taking into account the asymptotic behavior (2.88, 2.94), the ansatz (2.99) translates to the decompositions $`\mathrm{\Phi }(t,x,y)`$ $`=V(x)g(u,v)V^{}(x)\text{with}g(u,v)\text{U(2)},`$ (2.101) $`\mathrm{{\rm Y}}(t,x,y)`$ $`=V(x)\chi (u,v)V^{}(x)\text{with}\chi (u,v)u(2).`$ (2.102) To aim for the sine-Gordon equation, one imposes certain algebraic constraints on $`a`$ and $`b`$ (and therefore on $`\psi `$). Their precise form, however, is not needed, as we are ultimately interested only in $`g`$ or $`\chi `$. Therefore, we instead directly restrict $`g(u,v)`$ to the form $$g=\left(\begin{array}{cc}g_+& 0\\ 0& g_{}\end{array}\right)=g_+P_++g_{}P_{}\text{with}g_+\text{U(1)}_+\text{and}g_{}\text{U(1)}_{}$$ (2.103) and with projectors $`P_+=(\begin{array}{cc}1& 0\\ 0& 0\end{array})`$ and $`P_{}=(\begin{array}{cc}0& 0\\ 0& 1\end{array})`$. This imbeds $`g`$ into a U(1)$`\times `$U(1) subgroup of U(2). Note that $`g_+`$ and $`g_{}`$ do not commute, due to the implicit star product. Invoking the field redefinition (2.98) we infer that the corresponding reduction for $`\chi (u,v)`$ should be<sup>6</sup><sup>6</sup>6 Complex conjugates of scalar functions are denoted with a dagger to remind the reader of their noncommutativity. $$\chi =\mathrm{i}\left(\begin{array}{cc}0& h^{}\\ h& 0\end{array}\right)\text{with}h,$$ (2.104) with the “bridge relations” $`\alpha (hh^{})`$ $`=g_+^{}_ug_+=g_{}^{}_ug_{},`$ (2.105) $`\frac{1}{\alpha }_vh`$ $`=g_{}^{}g_+\mathrm{𝟏}\text{and h.c.}.`$ In this way, the $`u(2)`$-matrix $`\chi `$ is restricted to be off-diagonal. We now investigate in turn the consequences of the ansätze (2.101, 2.103) and (2.102, 2.104) for the equations of motion (2.90) and (2.96), respectively. ##### Reduction of Yang-type equation Let us insert the ansatz (2.101) into the Yang-type equation of motion (2.90). After stripping off the $`V`$ factors one obtains $$_v(g^{}_ug)+\alpha ^2(\sigma _1g^{}\sigma _1gg^{}\sigma _1g\sigma _1)=0.$$ (2.106) Specializing with (2.103) and employing the identities $`\sigma _1P_\pm \sigma _1=P_{}`$ we arrive at $`Y_+P_++Y_{}P_{}=0`$, with $`Y_+`$ $`_v(g_+^{}_ug_+)+\alpha ^2(g_{}^{}g_+g_+^{}g_{})=0,`$ (2.107) $`Y_{}`$ $`_v(g_{}^{}_ug_{})+\alpha ^2(g_+^{}g_{}g_{}^{}g_+)=0.`$ Since the brackets multiplying $`\alpha ^2`$ are equal and opposite, it is worthwhile to present the sum and the difference of the two equations: $`_v\left(g_+^{}_ug_++g_{}^{}_ug_{}\right)`$ $`=0,`$ (2.108) $`_v\left(g_+^{}_ug_+g_{}^{}_ug_{}\right)`$ $`=2\alpha ^2\left(g_+^{}g_{}g_{}^{}g_+\right).`$ It is natural to introduce the angle fields $`\varphi _\pm (u,v)`$ via $$g=e^{\frac{\mathrm{i}}{2}\varphi _+P_+}e^{\frac{\mathrm{i}}{2}\varphi _{}P_{}}g_+=e^{\frac{\mathrm{i}}{2}\varphi _+}\text{and}g_{}=e^{\frac{\mathrm{i}}{2}\varphi _{}}.$$ (2.109) In terms of these, the equations (2.108) read $`_v\left(e^{\frac{\mathrm{i}}{2}\varphi _+}_ue^{\frac{\mathrm{i}}{2}\varphi _+}+e^{\frac{\mathrm{i}}{2}\varphi _{}}_ue^{\frac{\mathrm{i}}{2}\varphi _{}}\right)`$ $`=0,`$ (2.110) $`_v\left(e^{\frac{\mathrm{i}}{2}\varphi _+}_ue^{\frac{\mathrm{i}}{2}\varphi _+}e^{\frac{\mathrm{i}}{2}\varphi _{}}_ue^{\frac{\mathrm{i}}{2}\varphi _{}}\right)`$ $`=2\alpha ^2\left(e^{\frac{\mathrm{i}}{2}\varphi _+}e^{\frac{\mathrm{i}}{2}\varphi _{}}e^{\frac{\mathrm{i}}{2}\varphi _{}}e^{\frac{\mathrm{i}}{2}\varphi _+}\right).`$ We propose to call these two equations “the noncommutative sine-Gordon equations”. Besides their integrability (see later sections for consequences) their form is quite convenient for studying the commutative limit. When $`\theta 0`$, (2.110) simplifies to $$_u_v(\varphi _+\varphi _{})=0\text{and}_u_v(\varphi _++\varphi _{})=8\alpha ^2\mathrm{sin}\frac{1}{2}(\varphi _++\varphi _{}).$$ (2.111) Because the equations have decoupled we may choose $$\varphi _+=\varphi _{}=:\varphi g_+=g_{}^{}g\text{U(1)}_\text{A}$$ (2.112) and reproduce the familiar sine-Gordon equation $$(_t^2_y^2)\varphi =4\alpha ^2\mathrm{sin}\varphi .$$ (2.113) One learns that in the commutative case the reduction is SU(2)$``$U(1)$`_\text{A}`$ since the U(1)$`_\text{V}`$ degree of freedom $`\varphi _+\varphi _{}`$ is not needed. The deformed situation, however, requires extending SU(2) to U(2), and so it is imperative here to keep both U(1)s and work with two scalar fields. Inspired by the commutative decoupling, one may choose another distinguished parametrization of $`g`$, namely $$g_+=e^{\frac{\mathrm{i}}{2}\rho }e^{\frac{\mathrm{i}}{2}\phi }\text{and}g_{}=e^{\frac{\mathrm{i}}{2}\rho }e^{\frac{\mathrm{i}}{2}\phi },$$ (2.114) which defines angles $`\rho (u,v)`$ and $`\phi (u,v)`$ for the linear combinations $`\text{U}(1)_\text{V}`$ and $`\text{U}(1)_\text{A}`$, respectively. Inserting this into (2.107) one finds $`_v\left(e^{\frac{\mathrm{i}}{2}\phi }_ue^{\frac{\mathrm{i}}{2}\phi }\right)+2\mathrm{i}\alpha ^2\mathrm{sin}\phi `$ $`=_v\left[e^{\frac{\mathrm{i}}{2}\phi }e^{\frac{\mathrm{i}}{2}\rho }(_ue^{\frac{\mathrm{i}}{2}\rho })e^{\frac{\mathrm{i}}{2}\phi }\right],`$ (2.115) $`_v\left(e^{\frac{\mathrm{i}}{2}\phi }_ue^{\frac{\mathrm{i}}{2}\phi }\right)2\mathrm{i}\alpha ^2\mathrm{sin}\phi `$ $`=_v\left[e^{\frac{\mathrm{i}}{2}\phi }e^{\frac{\mathrm{i}}{2}\rho }(_ue^{\frac{\mathrm{i}}{2}\rho })e^{\frac{\mathrm{i}}{2}\phi }\right].`$ In the commutative limit, this system is easily decoupled to $$_u_v\rho =0\text{and}_u_v\phi +4\alpha ^2\mathrm{sin}\phi =0,$$ (2.116) revealing that $`\rho \frac{1}{2}(\varphi _+\varphi _{})`$ and $`\phi \frac{1}{2}(\varphi _++\varphi _{})=\varphi `$ in this limit. It is not difficult to write down an action for (2.107) (and hence for (2.110) or (2.115)). The relevant action may be computed by reducing (2.92) with the help of (2.101) and (2.103). The result takes the form $$S[g_+,g_{}]=S_W[g_+]+S_W[g_{}]+\alpha ^2dtdy\left(g_+^{}g_{}+g_{}^{}g_+2\right),$$ (2.117) where $`S_W`$ is the abelian WZW action $$S_W[f]\frac{1}{2}dtdy_vf^1_uf\frac{1}{3}dtdy_0^1d\lambda ϵ^{\mu \nu \sigma }\widehat{f}^1_\mu \widehat{f}\widehat{f}^1_\nu \widehat{f}\widehat{f}^1_\sigma \widehat{f}$$ (2.118) Here $`\widehat{f}(\lambda )`$ is a homotopy path satisfying the conditions $`\widehat{f}(0)=1`$ and $`\widehat{f}(1)=f`$. Parametrizing $`g_\pm `$ as in (2.114) and using the Polyakov-Wiegmann identity, the action for $`\rho `$ and $`\phi `$ reads $`S[\rho ,\phi ]`$ $`=2S_{PC}\left[e^{\frac{\mathrm{i}}{2}\phi }\right]+\mathrm{\hspace{0.17em}2}\alpha ^2{\displaystyle dtdy\left(\mathrm{cos}\phi 1\right)}+\mathrm{\hspace{0.17em}2}S_W\left[e^{\frac{\mathrm{i}}{2}\rho }\right]`$ (2.119) $`{\displaystyle dtdye^{\frac{\mathrm{i}}{2}\rho }_ve^{\frac{\mathrm{i}}{2}\rho }\left(e^{\frac{\mathrm{i}}{2}\phi }_ue^{\frac{\mathrm{i}}{2}\phi }+e^{\frac{\mathrm{i}}{2}\phi }_ue^{\frac{\mathrm{i}}{2}\phi }\right)},`$ where $$S_{PC}[f]\frac{1}{2}dtdy_vf^1_uf.$$ (2.120) In this parametrization the WZ term has apparently been shifted entirely to the $`\rho `$ field while the cosine-type self-interaction remains for the $`\phi `$ field only. This fact has important consequences for the scattering amplitudes. It is well known that in ordinary commutative geometry the bosonization of $`N`$ free massless fermions in the fundamental representation of SU($`N`$) gives rise to a WZW model for a scalar field in SU($`N`$) plus a free scalar field associated with the U(1) invariance of the fermionic system. In the noncommutative case the bosonization of a single massless Dirac fermion produces a noncommutative U(1) WZW model , which becomes free only in the commutative limit. Moreover, the U(1) subgroup of U($`N`$) does no longer decouple , so that $`N`$ noncommuting free massless fermions are related to a noncommutative WZW model for a scalar in U($`N`$). On the other hand, giving a mass to the single Dirac fermion leads to a noncommutative cosine potential on the bosonized side . In contrast, the noncommutative sine-Gordon model we propose in this paper is of a more general form. The action (2.117) describes the propagation of a scalar field $`g`$ taking its value in U(1)$`\times `$U(1) $``$ U(2). Therefore, we expect it to be a bosonized version of two fermions in some representation of U(1)$`\times `$U(1). The absence of a WZ term for $`\phi `$ and the lack of a cosine-type self-interaction for $`\rho `$ as well as the non-standard interaction term make the precise identification non-trivial however. ##### Reduction of Leznov-type equation Alternatively, if we insert the ansatz (2.102) into the Leznov-type equation of motion (2.96) we get $$_u_v\chi +2\alpha ^2(\chi \sigma _1\chi \sigma _1)+\mathrm{i}\alpha [[\sigma _1,\chi ],_v\chi ]=0.$$ (2.121) Specializing with (2.104) this takes the form $`Z\sigma _{}+Z^{}\sigma _+=0`$ with $`\sigma _{}=(\begin{array}{cc}0& 0\\ 1& 0\end{array})`$ and $`\sigma _+=(\begin{array}{cc}0& 1\\ 0& 0\end{array})`$, where $$Z_u_vh+2\alpha ^2(hh^{})+\alpha \{_vh,hh^{}\}=0.$$ (2.122) The decomposition $$\chi =\mathrm{i}(h_1\sigma _1+h_2\sigma _2)h=h_1+\mathrm{i}h_2$$ (2.123) then yields $`_u_vh_12\alpha \{_vh_2,h_2\}`$ $`=0,`$ (2.124) $`_u_vh_2+4\alpha ^2h_2+2\alpha \{_vh_1,h_2\}`$ $`=0.`$ These two equations constitute an alternative description of the noncommutative sine-Gordon model; they are classically equivalent to the pair of (2.108) or, to be more specific, to the pair of (2.115). For the real fields the “bridge relations” (2.105) read $`2\mathrm{i}\alpha h_2=e^{\frac{\mathrm{i}}{2}\phi }e^{\frac{\mathrm{i}}{2}\rho }_u(e^{\frac{\mathrm{i}}{2}\rho }e^{\frac{\mathrm{i}}{2}\phi })=e^{\frac{\mathrm{i}}{2}\phi }e^{\frac{\mathrm{i}}{2}\rho }_u(e^{\frac{\mathrm{i}}{2}\rho }e^{\frac{\mathrm{i}}{2}\phi }),`$ (2.125) $`\frac{1}{\alpha }_vh_1=\mathrm{cos}\phi 1\text{and}\frac{1}{\alpha }_vh_2=\mathrm{sin}\phi .`$ One may “solve” one equation of (2.115) by an appropriate field redefinition from (2.125), which implies already one member of (2.124). The second equation from (2.115) then yields the remaining “bridge relations” in (2.125) as well as the other member of (2.124). This procedure works as well in the opposite direction, from (2.124) to (2.115). The nonlocal duality between $`(\phi ,\rho )`$ and $`(h_1,h_2)`$ is simply a consequence of the equivalence between (2.90) and (2.96) which in turn follows from our linear system (2.83). The “$`h`$ description” has the advantage of being polynomial. It is instructive to expose the action for the system (2.124). Either by inspection or by reducing the Leznov action (2.97) one obtains $$S[h_1,h_2]=dtdy\left\{_uh_1_vh_1+_uh_2_vh_24\alpha ^2h_2^24\alpha h_2^2_vh_1\right\}.$$ (2.126) #### 2.3.4 Relation with the previous noncommutative generalization of the sine-Gordon model The noncommutative generalizations of the sine-Gordon model presented above are expected to possess an infinite number of conservation laws, as they originate from the reduction of an integrable model . It is worthwhile to point out their relation to the noncommutative sine-Gordon model I discussed in section 2.2, which also features an infinite number of local conserved currents. In an alternative noncommutative version of the sine-Gordon model was proposed. Using the bicomplex approach the equations of motion were obtained as flatness conditions of a bidifferential calculus,<sup>7</sup><sup>7</sup>7 This subsection switches to Euclidean space $`^2`$, where $``$ and $`\overline{}`$ are derivatives with respect to complex coordinates. $$\overline{}(G^1G)=[R,G^1RG]_{},$$ (2.127) where $$R=2\alpha \left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)$$ (2.128) and $`G`$ is a suitable matrix in U(2) or, more generally, in complexified U(2). In the $`G`$ matrix was chosen as $$G=e_{}^{\frac{\mathrm{i}}{2}\sigma _2\mathrm{\Phi }}=\left(\begin{array}{cc}\mathrm{cos}_{}\frac{\mathrm{\Phi }}{2}& \mathrm{sin}_{}\frac{\mathrm{\Phi }}{2}\\ \mathrm{sin}_{}\frac{\mathrm{\Phi }}{2}& \mathrm{cos}_{}\frac{\mathrm{\Phi }}{2}\end{array}\right)$$ (2.129) with $`\mathrm{\Phi }`$ being a complex scalar field. This choice produces the noncommutative equations (all the products are $``$-products) $`\overline{}\left(e^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}e^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}+e^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}e^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}\right)=0,`$ $`\overline{}\left(e^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}e^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}e^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}e^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}\right)=4\mathrm{i}\alpha ^2\mathrm{sin}\mathrm{\Phi }.`$ (2.130) As shown in these equations (or a linear combination of them) can be obtained as a dimensional reduction of the equations of motion for noncommutative U(2) SDYM in 2+2 dimensions. The equations (2.130) can also be derived from an action which consists of the sum of two WZW actions augmented by a cosine potential, $$S[f,\overline{f}]=S[f]+S[\overline{f}]\text{with}S[f]S_W[f]\alpha ^2dtdy\left(f^2+f^22\right),$$ (2.131) with $`S_W[f]`$ given in (2.118) for $`fe^{\frac{\mathrm{i}}{2}\mathrm{\Phi }}`$ in complexified U(1). However, this action cannot be obtained from the SDYM action in 2+2 dimensions by performing the same field parametrization which led to (2.130). Comparing the actions (2.117) and (2.131) and considering $`f`$ and $`\overline{f}`$ as independent U(1) group valued fields we are tempted to formally identify $`fg_+`$ and $`\overline{f}g_{}`$. Doing this, we immediately realize that the two models differ in their interaction term which generalizes the cosine potential. While in (2.131) the fields $`f`$ and $`\overline{f}`$ show only self-interaction, the fields $`g_+`$ and $`g_{}`$ in (2.117) interact with each other. As we will see in section 2.3.6 this makes a big difference when evaluating the S-matrix elements. We close this section by observing that the equations of motion (2.110) can also be obtained directly in two dimensions by using the bicomplex approach described in . In fact, if instead of (2.129) we choose $$G=\left(\begin{array}{cc}e^{\frac{\mathrm{i}}{2}\varphi _+}+e^{\frac{\mathrm{i}}{2}\varphi _{}}& \mathrm{i}e^{\frac{\mathrm{i}}{2}\varphi _+}+\mathrm{i}e^{\frac{\mathrm{i}}{2}\varphi _{}}\\ \mathrm{i}e^{\frac{\mathrm{i}}{2}\varphi _+}\mathrm{i}e^{\frac{\mathrm{i}}{2}\varphi _{}}& e^{\frac{\mathrm{i}}{2}\varphi _+}+e^{\frac{\mathrm{i}}{2}\varphi _{}}\end{array}\right)$$ (2.132) it is easy to prove that (2.127) yields exactly the set of equations (2.110). Therefore, by exploiting the results in it should be straightforward to construct the first nontrivial conserved currents for the present model. #### 2.3.5 Solitons ##### Dressing approach in 2+1 dimensions. The existence of the linear system allows for powerful methods to systematically construct explicit solutions for $`\mathrm{\Psi }`$ and hence for $`\mathrm{\Phi }^{}=\mathrm{\Psi }\|_{\zeta =0}`$ or $`\mathrm{{\rm Y}}`$. For our purposes the so-called dressing method proves to be most practical, and so we shall first present it here for our linear system (2.83), before reducing the results to solitonic solutions of the noncommutative sine-Gordon equations. The central idea is to demand analyticity in the spectral parameter $`\zeta `$ for the linear system (2.83), which strongly restricts the possible form of $`\mathrm{\Psi }`$. The most elegant way to exploit this constraint starts from the observation that the left hand sides of the differential relations (D):=(2.84) as well as the reality condition (R):=(2.85) do not depend on $`\zeta `$ while their right hand sides are expected to be nontrivial functions of $`\zeta `$ (except for the trivial case $`\mathrm{\Psi }=\mathrm{\Psi }^0`$). More specifically, $`P^1`$ being compact, the matrix function $`\mathrm{\Psi }(\zeta )`$ cannot be holomorphic everywhere but must possess some poles, and hence the right hand sides of (D) and (R) should display these (and complex conjugate) poles as well. The resolution of this conundrum demands that the residues of the right hand sides at any would-be pole in $`\zeta `$ have to vanish. We are now going to evaluate these conditions. The dressing method builds a solution $`\mathrm{\Psi }_N(t,x,y,\zeta )`$ featuring $`N`$ simple poles at positions $`\mu _1,\mu _2,\mathrm{},\mu _N`$ by left-multiplying an $`(N1)`$-pole solution $`\mathrm{\Psi }_{N1}(t,x,y,\zeta )`$ with a single-pole factor of the form $`\left(1+\frac{\mu _N\overline{\mu }_N}{\zeta \mu _N}P_N(t,x,y)\right)`$, where the $`n\times n`$ matrix function $`P_N`$ is yet to be determined. In addition, we are free to right-multiply $`\mathrm{\Psi }_{N1}(t,x,y,\zeta )`$ with some constant unitary matrix $`\widehat{\mathrm{\Psi }}_N^0`$. Starting from $`\mathrm{\Psi }_0=\mathrm{𝟏}`$, the iteration $`\mathrm{\Psi }_0\mathrm{\Psi }_1\mathrm{}\mathrm{\Psi }_N`$ yields a multiplicative ansatz for $`\mathrm{\Psi }_N`$ which, via partial fraction decomposition, may be rewritten in an additive form (as a sum of simple pole terms). Let us trace this iterative procedure constructively. In accord with the outline above, the one-pole ansatz must read ($`\widehat{\mathrm{\Psi }}_1^0=:\mathrm{\Psi }_1^0`$) $$\mathrm{\Psi }_1=\left(\mathrm{𝟏}+\frac{\mu _1\overline{\mu }_1}{\zeta \mu _1}P_1\right)\mathrm{\Psi }_1^0=\left(\mathrm{𝟏}+\frac{\mathrm{\Lambda }_{11}S_1^{}}{\zeta \mu _1}\right)\mathrm{\Psi }_1^0$$ (2.133) with some $`n\times r_1`$ matrix functions $`\mathrm{\Lambda }_{11}`$ and $`S_1`$ for some $`1r_1<n`$. The normalization matrix $`\mathrm{\Psi }_1^0`$ is constant and unitary. It is quickly checked that $$\mathrm{res}_{\zeta =\overline{\mu }_1}(R)=0P_1^{}=P_1=P_1^2P_1=T_1(T_1^{}T_1)^1T_1^{},$$ (2.134) meaning that $`P_1`$ is a rank $`r_1`$ projector built from an $`n\times r_1`$ matrix function $`T_1`$. The columns of $`T_1`$ span the image of $`P_1`$ and obey $`P_1T_1=T_1`$. When using the second parametrization of $`\mathrm{\Psi }_1`$ in (2.133) one finds that $$\mathrm{res}_{\zeta =\overline{\mu }_1}(R)=0(\mathrm{𝟏}P_1)S_1\mathrm{\Lambda }_{11}^{}=0T_1=S_1$$ (2.135) modulo a freedom of normalization. Finally, the differential relations yield $$\mathrm{res}_{\zeta =\overline{\mu }_1}(D)=0(\mathrm{𝟏}P_1)\overline{L}_1^{A,B}(S_1\mathrm{\Lambda }_{11}^{})=0\overline{L}_1^{A,B}S_1=S_1\mathrm{\Gamma }_1^{A,B}$$ (2.136) for some $`r_1\times r_1`$ matrices $`\mathrm{\Gamma }_1^A`$ and $`\mathrm{\Gamma }_1^B`$, after having defined $$\overline{L}_i^A:=_u\overline{\mu }_i_x\text{and}\overline{L}_i^B:=\mu _i(_x\overline{\mu }_i_v)\text{for}i=1,2,\mathrm{},N.$$ (2.137) Because the $`\overline{L}_i^{A,B}`$ are linear differential operators it is easy to write down the general solution for (2.136): Introduce “co-moving coordinates” $$w_i:=x+\overline{\mu }_iu+\overline{\mu }_i^1v\overline{w}_i=x+\mu _iu+\mu _i^1v\text{for}i=1,2,\mathrm{},N$$ (2.138) so that on functions of $`(w_i,\overline{w}_i)`$ alone the $`\overline{L}_i^{A,B}`$ act as $$\overline{L}_i^A=\overline{L}_i^B=(\mu _i\overline{\mu }_i)\frac{}{\overline{w}_i}.$$ (2.139) Hence, (2.136) is solved by $`S_1(t,x,y)=\widehat{S}_1(w_1)e^{\overline{w}_1\mathrm{\Gamma }_1/(\mu _1\overline{\mu }_1)}`$ (2.140) for any $`w_1`$-holomorphic $`n\times r_1`$ matrix function $`\widehat{S}_1`$ (2.141) and $`\mathrm{\Gamma }_1^A=\mathrm{\Gamma }_1^B=:\mathrm{\Gamma }_1`$. Appearing to the right of $`\widehat{S}_1`$, the exponential factor is seen to drop out in the formation of $`P_1`$ via (2.134) and (2.135). Thus, no generality is lost by taking $`\mathrm{\Gamma }_1=0`$. We learn that any $`w_1`$-holomorphic $`n\times r_1`$ matrix $`T_1`$ is admissible to build a projector $`P_1`$ which then yields a solution $`\mathrm{\Psi }_1`$ (and thus $`\mathrm{\Phi }`$) via (2.133). Note that $`\mathrm{\Lambda }_{11}`$ need not be determined seperately but follows from our above result. It is not necessary to also consider the residues at $`\zeta =\mu _1`$ since their vanishing leads merely to the hermitian conjugated conditions. Let us proceed to the two-pole situation. The dressing ansatz takes the form ($`\mathrm{\Psi }_1^0\widehat{\mathrm{\Psi }}_2^0=:\mathrm{\Psi }_2^0`$) $$\mathrm{\Psi }_2=\left(\mathrm{𝟏}+\frac{\mu _2\overline{\mu }_2}{\zeta \mu _2}P_2\right)\left(\mathrm{𝟏}+\frac{\mu _1\overline{\mu }_1}{\zeta \mu _1}P_1\right)\mathrm{\Psi }_2^0=\left(\mathrm{𝟏}+\frac{\mathrm{\Lambda }_{21}S_1^{}}{\zeta \mu _1}+\frac{\mathrm{\Lambda }_{22}S_2^{}}{\zeta \mu _2}\right)\mathrm{\Psi }_2^0,$$ (2.142) where $`P_2`$ and $`S_2`$ are to be determined but $`P_1`$ and $`S_1`$ can be copied from above. Indeed, inspecting the residues of (R) and (D) at $`\zeta =\overline{\mu }_1`$ simply confirms that $$P_1=T_1(T_1^{}T_1)^1T_1^{}\text{and}T_1=S_1\text{with}S_1=\widehat{S}_1(w_1)$$ (2.143) is just carried over from the one-pole solution. Relations for $`P_2`$ and $`S_2`$ arise from $`\mathrm{res}_{\zeta =\overline{\mu }_2}(R)=0`$ $`(\mathrm{𝟏}P_2)P_2=0P_2=T_2(T_2^{}T_2)^1T_2^{},`$ (2.144) $`\mathrm{res}_{\zeta =\overline{\mu }_2}(R)=0`$ $`\mathrm{\Psi }_2(\overline{\mu }_2)S_2\mathrm{\Lambda }_{22}^{}=(\mathrm{𝟏}P_2)(1\frac{\mu _1\overline{\mu }_1}{\mu _1\overline{\mu }_2}P_1)S_2\mathrm{\Lambda }_{22}^{}=0,`$ (2.145) where the first equation makes use of the multiplicative form of the ansatz (2.142) while the second one exploits the additive version. We conclude that $`P_2`$ is again a hermitian projector (of some rank $`r_2`$) and thus built from an $`n\times r_2`$ matrix function $`T_2`$. Furthermore, (2.145) reveals that $`T_2`$ cannot be identified with $`S_2`$ this time, but we rather have $$T_2=\left(1\frac{\mu _1\overline{\mu }_1}{\mu _1\overline{\mu }_2}P_1\right)S_2$$ (2.146) instead. Finally, we consider $$\mathrm{res}_{\zeta =\overline{\mu }_2}(D)=0\mathrm{\Psi }_2(\overline{\mu }_2)\overline{L}_2^{A,B}(S_2\mathrm{\Lambda }_{22}^{})=0\overline{L}_2^{A,B}S_2=S_2\mathrm{\Gamma }_2^{A,B}$$ (2.147) which is solved by $$S_2(t,x,y)=\widehat{S}_2(w_2)e^{\overline{w}_2\mathrm{\Gamma }_2/(\mu _2\overline{\mu }_2)}$$ (2.148) for any $`w_2`$-holomorphic $`n\times r_2`$ matrix function $`\widehat{S}_2`$ and $`\mathrm{\Gamma }_2^A=\mathrm{\Gamma }_2^B=:\mathrm{\Gamma }_2`$. Once more, we are entitled to put $`\mathrm{\Gamma }_2=0`$. Hence, the second pole factor in (2.142) is constructed in the same way as the first one, except for the small complication (2.146). Again, $`\mathrm{\Lambda }_{21}`$ and $`\mathrm{\Lambda }_{22}`$ can be read off the result if needed. It is now clear how the iteration continues. After $`N`$ steps the final result reads $$\mathrm{\Psi }_N=\left\{\underset{\mathrm{}=0}{\overset{N1}{}}\left(\mathrm{𝟏}+\frac{\mu _N\mathrm{}\overline{\mu }_N\mathrm{}}{\zeta \mu _N\mathrm{}}P_N\mathrm{}\right)\right\}\mathrm{\Psi }_N^0=\left\{\mathrm{𝟏}+\underset{i=1}{\overset{N}{}}\frac{\mathrm{\Lambda }_{Ni}S_i^{}}{\zeta \mu _i}\right\}\mathrm{\Psi }_N^0,$$ (2.149) featuring hermitian rank $`r_i`$ projectors $`P_i`$ at $`i=1,2,\mathrm{},N`$, via $$P_i=T_i(T_i^{}T_i)^1T_i^{}\text{with}T_i=\{\underset{\mathrm{}=1}{\overset{i1}{}}(\mathrm{𝟏}\frac{\mu _i\mathrm{}\overline{\mu }_i\mathrm{}}{\mu _i\mathrm{}\overline{\mu }_i}P_i\mathrm{})\}S_i,$$ (2.150) where $$S_i(t,x,y)=\widehat{S}_i(w_i)$$ (2.151) for arbitrary $`w_i`$-holomorphic $`n\times r_i`$ matrix functions $`\widehat{S}_i(w_i)`$. The corresponding classical Yang and Leznov fields are $`\mathrm{\Phi }_N`$ $`=\mathrm{\Psi }_N^{}(\zeta =0)=\mathrm{\Psi }_{N}^{0}{}_{}{}^{}{\displaystyle \underset{i=1}{\overset{N}{}}}\left(\mathrm{𝟏}\rho _iP_i\right)\text{with}\rho _i=1{\displaystyle \frac{\mu _i}{\overline{\mu }_i}},`$ (2.152) $`\mathrm{{\rm Y}}_N`$ $`=\underset{\zeta \mathrm{}}{lim}\zeta \left(\mathrm{\Psi }_N(\zeta )\mathrm{\Psi }_{N}^{0}{}_{}{}^{}\mathrm{𝟏}\right)={\displaystyle \underset{i=1}{\overset{N}{}}}(\mu _i\overline{\mu }_i)P_i.`$ (2.153) The solution space constructed here is parametrized (slightly redundantly) by the set $`\{\widehat{S}_i\}_1^N`$ of matrix-valued holomorphic functions and the pole positions $`\mu _i`$. The so-constructed classical configurations have solitonic character (meaning finite energy) when all these functions are algebraic. The dressing technique as presented above is well known in the commutative theory; novel is only the realization that it carries over verbatim to the noncommutative situation by simply understanding all products as star products (and likewise inverses, exponentials, etc.). Of course, it may be technically difficult to $``$-invert some matrix, but one may always fall back on an expansion in powers of $`\theta `$. ##### Solitons of the noncommutative sine-Gordon theory We should now be able to generate $`N`$-soliton solutions to the noncommutative sine-Gordon equations, say (2.115), by applying the reduction from $`2+1`$ to $`1+1`$ dimensions (see previous section) to the above strategy for the group U(2), i.e. putting $`n=2`$. In order to find nontrivial solutions, we specify the constant matrix $``$ in the ansatz (2.99) for $`\mathrm{\Psi }`$ as $$=e^{\mathrm{i}\frac{\pi }{4}\sigma _2}=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)$$ (2.154) which obeys the relations $`\sigma _3=\sigma _1`$ and $`\sigma _1=\sigma _3`$. Pushing $``$ beyond $`V`$ we can write $$\mathrm{\Phi }(t,x,y)=W(x)\stackrel{~}{g}(u,v)W^{}(x)\text{with}W(x)=e^{\mathrm{i}\alpha x\sigma _3}$$ (2.155) and $$\stackrel{~}{g}(u,v)=g(u,v)^{}=\left(\begin{array}{cc}g_+& 0\\ 0& g_{}\end{array}\right)^{}=\frac{1}{2}\left(\begin{array}{cc}g_++g_{}& g_+g_{}\\ g_+g_{}& g_++g_{}\end{array}\right).$$ (2.156) With hindsight from the commutative case we choose $$\widehat{\mathrm{\Psi }}_i^0=\sigma _3i\mathrm{\Psi }_N^0=\sigma _3^N$$ (2.157) (which commutes with $`W`$) and restrict the poles of $`\mathrm{\Psi }`$ to the imaginary axis, $`\mu _i=\mathrm{i}p_i`$ with $`p_i`$. Therewith, the co-moving coordinates (2.138) become $$w_i=x\mathrm{i}(p_iup_i^1v)=:x\mathrm{i}\eta _i(u,v),$$ (2.158) defining $`\eta _i`$ as real linear functions of the light-cone coordinates. Consequentially, from (2.152) we get $`\rho _i=2`$ and find that $$\stackrel{~}{g}_N(u,v)=\sigma _3^N\underset{i=1}{\overset{N}{}}\left(\mathrm{𝟏}2\stackrel{~}{P}_i(u,v)\right)\text{with}P_i=W\stackrel{~}{P}_iW^{}.$$ (2.159) Repeating the analysis of the previous subsection, one is again led to construct hermitian projectors $$\stackrel{~}{P}_i=\stackrel{~}{T}_i(\stackrel{~}{T}_i^{}\stackrel{~}{T}_i)^1\stackrel{~}{T}_i^{}\text{with}\stackrel{~}{T}_i=\underset{\mathrm{}=1}{\overset{i1}{}}\left(\mathrm{𝟏}\frac{2p_i\mathrm{}}{p_i\mathrm{}+p_i}\stackrel{~}{P}_i\mathrm{}\right)\stackrel{~}{S}_i,$$ (2.160) where $`2\times 1`$ matrix functions $`\stackrel{~}{S}_i(u,v)`$ are subject to $$\stackrel{~}{\overline{L}}_i^{A,B}\stackrel{~}{S}_i=\stackrel{~}{S}_i\stackrel{~}{\mathrm{\Gamma }}_i\text{for}i=1,2,\mathrm{},N$$ (2.161) and some numbers $`\stackrel{~}{\mathrm{\Gamma }}_i`$ (note that now rank $`r_i=1`$) which again we can put to zero. On functions of the reduced co-moving coordinates $`\eta _i`$ alone, $$\stackrel{~}{\overline{L}}_i^{A,B}=W^{}\overline{L}_i^{A,B}W=(\mu _i\overline{\mu }_i)W^{}\frac{}{\overline{w}_i}W=p_i\left(\frac{}{\eta _i}+\alpha \sigma _3\right)$$ (2.162) so that (2.161) is solved by $$\stackrel{~}{S}_i(u,v)=\widehat{\stackrel{~}{S}}_i(\eta _i)=\left(\begin{array}{c}\gamma _{i1}e^{\alpha \eta _i}\\ \mathrm{i}\gamma _{i2}e^{+\alpha \eta _i}\end{array}\right)=e^{\alpha \eta _i\sigma _3}\left(\begin{array}{c}\gamma _{i1}\\ \mathrm{i}\gamma _{i2}\end{array}\right)\text{with}\gamma _{i1},\gamma _{i2}.$$ (2.163) Furthermore, it is useful to rewrite $$\gamma _{i1}\gamma _{i2}=:\lambda _i^2\text{and}\gamma _{i2}/\gamma _{i1}=:\gamma _i^2\left(\begin{array}{c}\gamma _{i1}\\ \mathrm{i}\gamma _{i2}\end{array}\right)=\lambda _i\left(\begin{array}{c}\gamma _i^1\\ \mathrm{i}\gamma _i\end{array}\right)$$ (2.164) because then $`\|\gamma _i\|`$ may be absorbed into $`\eta _i`$ by shifting $`\alpha \eta _i\alpha \eta _i+\mathrm{ln}\|\gamma _i\|`$. The multipliers $`\lambda _i`$ drop out in the computation of $`\stackrel{~}{P}_i`$. Finally, to make contact with the form (2.156) we restrict the constants $`\gamma _i`$ to be real. Let us check the one-soliton solution, i.e. put $`N=1`$. Suppressing the indices momentarily, absorbing $`\gamma `$ into $`\eta `$ and dropping $`\lambda `$, we infer that $`\stackrel{~}{T}=\left(\begin{array}{c}e^{\alpha \eta }\\ \mathrm{i}e^{\alpha \eta }\end{array}\right)\stackrel{~}{P}={\displaystyle \frac{1}{2\mathrm{ch2}\alpha \eta }}\left(\begin{array}{cc}e^{2\alpha \eta }& \mathrm{i}\\ \mathrm{i}& e^{+2\alpha \eta }\end{array}\right)`$ (2.165) $`\stackrel{~}{g}=\left(\begin{array}{cc}\mathrm{th2}\alpha \eta & \frac{\mathrm{i}}{\mathrm{ch2}\alpha \eta }\\ \frac{\mathrm{i}}{\mathrm{ch2}\alpha \eta }& \mathrm{th2}\alpha \eta \end{array}\right)`$ (2.166) which has $`\mathrm{det}\stackrel{~}{g}=1`$. Since here the entire coordinate dependence comes in the single combination $`\eta (u,v)`$, all star products trivialize and the one-soliton configuration coincides with the commutative one. Hence, the field $`\rho `$ drops out, $`\stackrel{~}{g}`$ SU(2), and we find, comparing (2.165) with (2.156), that $$\frac{1}{2}(g_++g_{})=\mathrm{cos}\frac{\phi }{2}=\mathrm{th2}\alpha \eta \text{and}\frac{1}{2\mathrm{i}}(\mathrm{g}_+\mathrm{g}_{})=\mathrm{sin}\frac{\phi }{2}=\frac{1}{\mathrm{ch2}\alpha \eta }$$ (2.167) which implies $$\mathrm{tan}\frac{\phi }{4}=e^{2\alpha \eta }\phi =4\mathrm{arctan}e^{2\alpha \eta }=2\mathrm{arcsin}(\mathrm{th2}\alpha \eta ),$$ (2.168) reproducing the well known sine-Gordon soliton with mass $`m=2\alpha `$. Its moduli parameters are the velocity $`\nu =\frac{1p^2}{1+p^2}`$ and the center of inertia $`y_0=\frac{1}{\alpha }\sqrt{1\nu ^2}\mathrm{ln}\|\gamma \|`$ at zero time . In passing we note that in the “$`h`$ description” the soliton solution takes the form $$h_1=p\mathrm{th2}\alpha \eta \text{and}\mathrm{h}_2=\frac{\mathrm{p}}{\mathrm{ch2}\alpha \eta }\mathrm{h}=\mathrm{p}\mathrm{th}(\alpha \eta +\frac{\mathrm{i}\pi }{4})=\mathrm{p}\mathrm{e}^{\frac{\mathrm{i}}{2}\phi }.$$ (2.169) Noncommutativity becomes relevant for multi-solitons. At $`N=2`$, for instance, one has $`\stackrel{~}{g}_2=(12\stackrel{~}{P}_1)(12\stackrel{~}{P}_2)\text{with}\stackrel{~}{P}_1=\stackrel{~}{P}\text{from (}\text{2.165}\text{)}`$ (2.170) $`\text{and}\stackrel{~}{P}_2=\stackrel{~}{T}_2(\stackrel{~}{T}_2^{}\stackrel{~}{T}_2)^1\stackrel{~}{T}_2^{}`$ $`\text{where}\stackrel{~}{T}_2=\left(\mathrm{𝟏}\frac{2p_1}{p_1+p_2}\stackrel{~}{P}_1\right)\widehat{\stackrel{~}{S}}_2\text{and}\widehat{\stackrel{~}{S}}_2=e^{\alpha \eta _2\sigma _3}\left(\begin{array}{c}\gamma _2^1\\ \mathrm{i}\gamma _2\end{array}\right)`$ $`\text{with}\gamma _2.`$ We refrain from writing down the lengthy explicit expression for $`\stackrel{~}{g}_2`$ in terms of the noncommuting coordinates $`\eta _1`$ and $`\eta _2`$, but one cannot expect to find a unit (star-)determinant for $`\stackrel{~}{g}_2`$ except in the commutative limit. This underscores the necessity of extending the matrices to U(2) and the inclusion of a nontrivial $`\rho `$ at the multi-soliton level. It is not surprising that the just-constructed noncommutative sine-Gordon solitons themselves descend directly from BPS solutions of the $`2+1`$ dimensional integrable sigma model. Indeed, putting back the $`x`$ dependence via (2.155), the $`2+1`$ dimensional projectors $`P_i`$ are built from $`2\times 1`$ matrices $$S_i=W(x)\widehat{S}_i(\eta _i)=ϵ^{\mathrm{i}\alpha w_i\sigma _3}\left(\begin{array}{c}\gamma _i^1\\ \mathrm{i}\gamma _i\end{array}\right)=\left(\begin{array}{c}1\\ \mathrm{i}\gamma _i^2ϵ^{2\mathrm{i}\alpha w_i}\end{array}\right)\gamma _i^1e^{\mathrm{i}\alpha w_i}.$$ (2.171) In the last expression the right factor drops out on the computation of projectors; the remaining column vector agrees with the standard conventions . Reassuringly, the coordinate dependence has combined into $`w_i`$. The ensueing $`2+1`$ dimensional configurations $`\mathrm{\Phi }_N`$ are nothing but noncommutative multi-plane-waves the simplest examples of which were already investigated in . #### 2.3.6 (Nice) properties of the S-matrix In this section we compute tree-level amplitudes for the noncommutative generalization of the sine-Gordon model proposed in section 2.3.3, both in the Yang and the Leznov formulation. In commutative geometry the sine-Gordon S-matrix factorizes in two-particle processes and no particle production occurs, as a consequence of the existence of an infinite number of conservation laws. In the noncommutative case it is interesting to investigate whether the presence of an infinite number of conserved currents is still sufficient to guarantee the integrability of the system in the sense of having a factorized S-matrix. The previous noncommutative version of the sine-Gordon model we introduced and studied in section 2.2 is endowed with an infinite set of conserved currents. In section 2.2.9 we have seen that, despite the existence of an infinite chain of conservation laws, particle production occurs in that model and that the S-matrix is neither factorized nor causal.<sup>8</sup><sup>8</sup>8 Acausal behaviour in noncommutative field theory was first observed in and shown to be related to time-space noncommutativity. As already stressed in section 2.3.4, the noncommutative generalization of the sine-Gordon model we proposed in and discussed in the present section 2.3 differs from the one studied in in the generalization of the cosine potential. Therefore, both theories describe the dynamics of two real scalar fields, but the structure of the interaction terms between the two fields is different. We then expect the scattering amplitudes of the present theory to behave differently from those of the previous one. To this end we will compute the amplitudes corresponding to $`22`$ processes for the fields $`\rho `$ and $`\phi `$ in the $`g`$-model (Yang formulation) as well as for the fields $`h_1`$ and $`h_2`$ in the $`h`$-model (Leznov formulation). In the $`g`$-model we will also compute $`24`$ and $`33`$ amplitudes for the massive field $`\phi `$. In both models the S-matrix will turn out to be factorized and causal in spite of their time-space noncommutativity. ##### Amplitudes in the “$`g`$-model”. Feynman rules We parametrize the $`g`$-model with $`(\rho ,\phi )`$ as in (2.119) since in this parametrization the mass matrix turns out to be diagonal, with zero mass for $`\rho `$ and $`m=2\alpha `$ for $`\phi `$. Expanding the action (2.119) up to the fourth order in the fields, we read off the following Feynman rules: * The propagators $``$ $`\phi \phi ={\displaystyle \frac{2\mathrm{i}}{k^24\alpha ^2}},`$ (2.172) $``$ $`\rho \rho ={\displaystyle \frac{2\mathrm{i}}{k^2}}.`$ (2.173) * The vertices (including a factor of “i” from the expansion of $`ϵ^{\mathrm{i}S}`$) $`=`$ $`{\displaystyle \frac{1}{2^3}}(k_2^2k_1^22k_1k_2)F(k_1,k_2,k_3)`$ (2.174) $`=`$ $`{\displaystyle \frac{1}{23!}}k_1k_2F(k_1,k_2,k_3)`$ (2.175) $`=`$ $`\left[{\displaystyle \frac{\mathrm{i}}{2^34!}}(k_1^2+3k_1k_3)+{\displaystyle \frac{2\mathrm{i}\alpha ^2}{4!}}\right]F(k_1,k_2,k_3,k_4)`$ (2.177) $`=`$ $`{\displaystyle \frac{\mathrm{i}}{2^34!}}(k_1^2+3k_1k_3)F(k_1,k_2,k_3,k_4)`$ (2.179) $`=`$ $`{\displaystyle \frac{\mathrm{i}}{2^5}}(k_1^2k_2^2+2k_1k_32k_2k_3+\mathrm{\hspace{0.17em}2}k_1k_2+2k_1k_3`$ (2.180) $`+2k_3k_2)F(k_1,k_2,k_3,k_4)`$ where we used the conventions of section 2 with the definitions $$uv=\eta ^{ab}u_av_b=u_tv_tu_yv_y\text{and}uv=u_tv_yu_yv_t.$$ (2.181) Moreover, we have defined $$F(k_1,\mathrm{},k_n)=\mathrm{exp}\left\{\frac{\mathrm{i}}{2}\theta _{i<j}^nk_ik_j\right\}.$$ (2.182) and use the convention that all momentum lines are entering the vertex and energy-momentum conservation has been taken into account. We now compute the scattering amplitudes $`\phi \phi \phi \phi `$, $`\rho \rho \rho \rho `$ and $`\phi \rho \phi \rho `$ and the production amplitude $`\phi \phi \rho \rho `$. We perform the calculations in the center-of-mass frame. We assign the convention that particles with momenta $`k_1`$ and $`k_2`$ are incoming, while those with momenta $`k_3`$ and $`k_4`$ are outgoing. ##### Amplitude $`\phi \phi \phi \phi `$ The four momenta are explicitly written as $`k_1=(E,p),k_2=(E,p),k_3=(E,p),k_4=(E,p),`$ (2.183) with the on-shell condition $`E^2p^2=4\alpha ^2`$. There are two topologies of diagrams contributing to this process. Taking into account the leg permutations corresponding to the same particle at a single vertex, the contributions read \| \| = \| $`2\mathrm{i}\alpha ^2\mathrm{cos}^2(\theta Ep),`$ \| \| = \| $`0,`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| \| = \| $`\frac{\mathrm{i}}{2}p^2\mathrm{sin}^2(\theta Ep),`$ \| \| = \| $`\frac{\mathrm{i}}{2}E^2\mathrm{sin}^2(\theta Ep).`$ \| The second diagram is actually affected by a collinear divergence since the total momentum $`k_1+k_4`$ for the internal massless particle is on-shell vanishing. We regularize this divergence by temporarily giving a small mass to the $`\rho `$ particle. It is easy to see that the amplitude is zero for any value of the small mass since the wedge products $`k_1k_4`$ and $`k_2k_3`$ from the two vertices always vanish. As an alternative procedure we can put one of the external particles slightly off-shell, so obtaining a finite result which vanishes in the on-shell limit. Summing all the contributions, for the $`\phi \phi \phi \phi `$ amplitude we arrive at $$A_{\phi \phi \phi \phi }=2\mathrm{i}\alpha ^2,$$ (2.184) which perfectly describes a causal amplitude. A nonvanishing $`\phi \phi \phi \phi `$ amplitude appears also in the noncommutative sine-Gordon proposal of . However, there the amplitude has a nontrivial $`\theta `$-dependence which is responsible for acausal behavior. Comparing the present result with the result in , we observe that the same kind of diagrams contribute. The main difference is that the exchanged particle is now massless instead of massive. This crucial difference leads to the cancellation of the $`\theta `$-dependent trigonometric behaviour which in the previous case gave rise to acausality. ##### Amplitude $`\rho \rho \rho \rho `$ In this case the center-of-mass momenta are given by $$k_1=(E,E),k_2=(E,E),k_3=(E,E),k_4=(E,E),$$ (2.185) where the on-shell condition $`E^2p^2=0`$ has already been taken into account. For this amplitude we have the following contributions \| \| = \| $`0,`$ \| \| = \| $`0,`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| \| = \| $`\frac{\mathrm{i}}{2}E^2\mathrm{sin}^2(\theta E^2),`$ \| \| = \| $`\frac{\mathrm{i}}{2}E^2\mathrm{sin}^2(\theta E^2).`$ \| Again, a collinear divergence appears in the second diagram. In order to regularize the divergence we can proceed as before by assigning a small mass to the $`\rho `$ particle. The main difference with respect to the previous case is that now the $`\rho `$ particle also appears as an external particle, with the consequence that the on-shell momenta in (2.185) will get modified by the introduction of a regulator mass. A careful calculation shows that the amplitude is zero for any value of the regulator mass, due to the vanishing of the factors $`k_1k_4`$ and $`k_2k_3`$ from the vertices. Therefore, the two nonvanishing contributions add to $$A_{\rho \rho \rho \rho }=0.$$ (2.186) ##### Amplitude $`\phi \rho \phi \rho `$ There are two possible configurations of momenta in the center-of-mass frame, describing the scattering of the massive particle with either a left-moving or a right-moving massless one. In the left-moving case the momenta are $`k_1=(E,p),k_2=(p,p),k_3=(E,p),k_4=(p,p),`$ (2.187) while in the right-moving case we have $`k_1=(E,p),k_2=(p,p),k_3=(E,p),k_4=(p,p).`$ (2.188) For the left-moving case (2.187) the results are \| \| = \| $`\frac{\mathrm{i}}{2}Ep\mathrm{sin}(\theta Ep)\mathrm{sin}(\theta p^2),`$ \| \| \| \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| \| = \| $`\frac{\mathrm{i}}{2}Ep\mathrm{sin}(\theta Ep)\mathrm{sin}(\theta p^2),`$ \| \| \| \| \| \| = \| $`0,`$ \| \| = \| $`0.`$ \| For the right-moving choice (2.188), we obtain instead \| \| = \| $`0,`$ \| \| = \| $`0,`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| \| = \| $`0,`$ \| \| = \| $`0.`$ \| In this second case an infrared divergence is present due to the massless propagator, but again it can be cured as described before. In both cases the scattering amplitude vanishes, $$A_{\phi \rho \phi \rho }=0.$$ (2.189) ##### Amplitude $`\phi \phi \rho \rho `$ The momenta in the center-of-mass frame are given by $`k_1=(E,p),k_2=(E,p),k_3=(E,E),k_4=(E,E).`$ (2.190) In this case we have three kinds of diagrams contributing. The corresponding results are \| \| = \| $`\frac{\mathrm{i}}{2}Ep\mathrm{sin}(\theta Ep)\mathrm{sin}(\theta E^2),`$ \| \| \| \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| \| = \| $`\frac{\mathrm{i}}{2}Ep\mathrm{sin}(\theta Ep)\mathrm{sin}(\theta E^2),`$ \| \| \| \| \| \| = \| $`0,`$ \| \| = \| $`0.`$ \| Summing the four contributions, we obtain $$A_{\phi \phi \rho \rho }=0$$ (2.191) as it should be expected for a production amplitude in an integrable model. The same is true for the time-reversed production, $$A_{\rho \rho \phi \phi }=0.$$ (2.192) Summarizing, we have found that the only nonzero amplitude for tree-level $`22`$ processes is the one describing the scattering among two of the massive excitations. The result is constant, independent of the momenta and so describes a perfectly causal process. Since the result is independent of the noncommutation parameter $`\theta `$ it agrees with the four-point amplitude for the ordinary sine-Gordon model. Finally, we have found that the production amplitudes $`\phi \phi \rho \rho `$ and $`\rho \rho \phi \phi `$ vanish, as required for ordinary integrable theories. As a further check of our calculation and an additional test of our model we have computed the production amplitude $`\phi \phi \phi \phi \phi \phi `$ and the scattering amplitude $`\phi \phi \phi \phi \phi \phi `$. In both cases the topologies we have to consider are \| \| \| \| \| --- \| --- \| --- \| \| \| \| \| . \| \| --- \| --- \| --- \| --- \| Due to the growing number of channels and ordering of vertices, it is no longer practical to perform the calculations by hand. We have used Mathematica<sup>©</sup> to symmetrize the vertices and take automatically into account the different diagrams obtained by exchanging momenta entering a given vertex. The computation has been performed with assigned values of the external momenta but arbitrary values for $`\alpha ^2`$ and $`\theta `$. We have found a vanishing result for both the scattering and the production amplitude. This is in agreement with the commutative sine-Gordon model results. ##### Amplitudes in the “$`h`$-model” We now discuss the $`22`$ amplitudes in the Leznov formulation. The theory is again described by two interacting fields, $`h_1`$ massless and $`h_2`$ massive. Referring to the action (2.126) we extract the following Feynman rules, * The propagators $``$ $`h_1h_1={\displaystyle \frac{\mathrm{i}}{2k^2}},`$ (2.193) $``$ $`h_2h_2={\displaystyle \frac{\mathrm{i}/2}{k^24\alpha ^2}}.`$ (2.194) * The vertex $`=`$ $`4\alpha (k_{3t}k_{3y})F(k_1,k_2,k_3).`$ (2.195) Again, we compute scattering amplitudes in the center-of-mass frame. Given the particular structure of the vertex, at tree level there is no $`h_1h_1h_1h_1`$ scattering. To find the $`h_2h_2h_2h_2`$ amplitude we assign the momenta (2.183) to the external particles. The contributions are \| \| = \| $`16\mathrm{i}\alpha ^2\mathrm{cos}^2(\theta Ep),`$ \| \| = \| $`16\mathrm{i}\alpha ^2\mathrm{cos}^2(\theta Ep),`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| \| = \| $`0.`$ \| \| \| \| We note that a collinear divergence appears in the last diagram which can be regularized as described before. Summing the two nonvanishing contributions we obtain complete cancellation. For the $`h_2h_2h_1h_1`$ amplitude the center-of-mass-momenta are given in (2.190). The only topology contributing to this production amplitude has two channels, yielding \| \| = \| $`0,`$ \| \| = \| $`0,`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| which are both zero, so giving a vanishing result once more. The same is true for the $`h_1h_1h_2h_2`$ production process. Finally, for the $`h_1h_2h_1h_2`$ amplitude, we refer to the center-of-mass momenta defined in (2.187) and (2.188). In both cases the contributions are \| \| = \| $`0,`$ \| \| = \| $`0,`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| and so we find that the sum of the two channels is always equal to zero. Since all the $`22`$ amplitudes vanish, the S-matrix is trivially causal and factorized. Both in the ordinary and noncommutative cases the “$`h`$-model” is dual to the “$`g`$-model”. In the commutative limit the “$`g`$-model” gives rise to a sine-Gordon model plus a free field which can be set to zero. In this limit our amplitudes exactly reproduce the sine-Gordon amplitudes. On the other hand, the amplitudes for the “$`h`$-model” all vanish. Therefore, in the commutative limit they do not reproduce anything immediately recognizable as an ordinary sine-Gordon amplitude. This can be understood by observing that, both in the ordinary and in the noncommutative case, the Leznov formulation is an alternative description of the sine-Gordon dynamics and obtained from the standard Yang formulation by the nonlocal field redefinition given in (2.125). Therefore, it is expected that the scattering amplitudes for the elementary exitations, which are different in the two formulations, do not resemble each other. #### 2.3.7 Conclusions In this section 2.3 I have introduced and discussed a novel noncommutative sine-Gordon system based on two scalar fields, which seems to retain all advantages of $`1+1`$ dimensional integrable models known from the commutative limit. The rationale for introducing a second scalar field was provided by deriving the sine-Gordon equations and action through dimensional and algebraic reduction of an integrable $`2+1`$ dimensional sigma model: In the noncommutative extension of this scheme it is natural to generalize the algebraic reduction of SU(2)$``$U(1) to one of U(2)$``$U(1)$`\times `$U(1). We gave two Yang-type and one Leznov-type parametrizations of the coupled system in (2.110), (2.115) and (2.124) and provided the actions for them, including a comparison with previous proposals. It was then outlined how to explicitly construct noncommutative sine-Gordon multi-solitons via the dressing method based on the underlying linear system. We found that the one-soliton configuration agrees with the commutative one but already the two-soliton solutions gets Moyal deformed. What is the gain of doubling the field content as compared to the standard sine-Gordon system or its straightforward star deformation? Usually, time-space noncommutativity adversely affects the causality and unitarity of the S-matrix (see, e.g. ), even in the presence of an infinite number of local conservation laws. In contrast, the model described in seems to possess an S-matrix which is causal and factorized, as we checked for all tree-level $`22`$ processes both in the Yang and Leznov formulations. Furthermore, we verified the vanishing of some $`33`$ scattering amplitudes and $`24`$ production amplitudes thus proving the absence of particle production. It would be nice to understand what actually drives a system to be integrable in the noncommutative case. A hint in this direction might be that the model proposed in has been constructed directly in two dimensions even if its equations of motion (but not the action) can be obtained by a suitable reduction of a four dimensional system (noncommutative self-dual Yang-Mills). The model proposed in this paper, instead, originates directly, already at the level of the action, from the reduction of noncommutative self-dual Yang-Mills theory which is known to be integrable and related to the $`N=2`$ string . Several directions of future research are suggested by our results. First, one might hope that our noncommutative two-field sine-Gordon model is equivalent to some two-fermion model via noncommutative bosonization. Second, it would be illuminating to derive the exact two-soliton solution and extract its scattering properties, either directly in our model or by reducing wave-like solutions of the 2+1 dimensional sigma model . Third, there is no obstruction against applying the ideas and techniques of this paper to other 1+1 dimensional noncommutative integrable systems in order to cure their pathologies as well. ## Chapter 3 Covariant superstring vertices and a possible nonconstant superspace deformation ### 3.1 An introduction to the pure spinor superstring In this section I will mostly refer to the review paper . #### 3.1.1 Motivation: Problems with RNS and GS formalisms In section 1.2.2 I have briefly outlined how the discussion of the stringy origin of bosonic noncommutative geometry presented in section 1.2.1 can be generalized to the superstring, in both RNS and GS formalisms. As anticipated, both of them display some awkward features, due to their target-space or worldsheet symmetry structure, respectively. The RNS formalism is characterized by an $`N=1`$ worldsheet supersymmetry. The field content is the set of bosonic coordinate fields $`x^m`$ (worldsheet scalars and spacetime vectors) together with the worldsheet spinors (and spacetime vectors) $`\psi ^m`$. The worldsheet action for the string in a flat background is quadratic, therefore the quantization in this formalism is straightforward. After a suitable consistent truncation of the spectrum (GSO projection), the theory also enjoys target space supersymmetry, but clearly this symmetry is not manifest. As a result, a series of problems arises, for instance in the computation of amplitudes with more than four external fermions and in dealing with general R-R backgrounds. The GS formalism instead is manifestly target-space supersymmetric, but the worldsheet symmetry structure is quite complicated. Target space is a ten-dimensional superspace described by the bosonic coordinates $`x^m`$ and their superpartners $`\theta ^\alpha `$, $`\widehat{\theta }^{\widehat{\beta }}`$ (in type II case). For the number of physical fermionic degrees of freedom to be related to the bosonic ones as required by target-space supersymmetry, a worldsheet fermionic local symmetry must be present ($`\kappa `$-symmetry, ). Therefore, the natural supersymmetric generalization of the bosonic string action $$S_1=d^2\sigma \sqrt{h}h^{ij}\mathrm{\Pi }_i\mathrm{\Pi }_j$$ (3.1) where $`h^{ij}`$ is the worldsheet metric and $`\mathrm{\Pi }_i^m`$ is the supersymmetrized bosonic momentum, does not work, not being $`\kappa `$-symmetric. When $`N2`$ a WZ term $`S_2`$ can be added, so that the resulting action $`S=S_1+S_2`$ is $`\kappa `$-symmetric (for this discussion the reader can refer to and references therein). $`S`$ in conformal gauge and in a flat background is given by (1.217) with $`B=0`$. It is nonquadratic and describes a complicated, interacting worldsheet field theory. This fact prevents quantization except in light-cone, where the action reduces to a quadratic form. Since light-cone gauge is not manifestly Lorentz covariant, problems in the computation of amplitudes emerge also in this formalism and only four-point tree and one-loop amplitudes have been computed. Moreover, backgrounds that do not allow for a light-cone choice cannot be dealt with at the quantum level. An alternative approach to the GS formalism was introduced by Siegel . The main problem of the GS superstring is that a set of phase-space constraints arise at the classical level whose structure do not allow for a Dirac quantization procedure. In particular, since the conjugate momenta $`p_\alpha `$ to the fermionic variables $`\theta ^\alpha `$ do not appear in the action, one has phase space constraints $`d_\alpha =0`$ together with the Virasoro constraint $`T=\frac{1}{2}\mathrm{\Pi }\mathrm{\Pi }=0`$ related to the conformal gauge choice. The anticommutator of the $`d`$’s is proportional to the $`\mathrm{\Pi }`$’s. As a result, half of the fermionic constraints are first class and half are second class . The separation of the two different kinds of constraint cannot be achieved in a manifestly Lorentz covariant way. This explains why quantization of the model only works in light-cone gauge. In Siegel proposed to rewrite the GS action in a first order formalism for the fermionic variables, hoping that a set of phase space constraints that are all first class could be found. These contraints were to be constructed out of the supersymmetric objects $`\mathrm{\Pi }^m`$, $`\theta ^\alpha `$ and the GS constraint $`d_\alpha `$, no longer constrained to vanish. The explicit form of the left-moving contribution to the GS action in conformal gauge in Siegel formalism is $$S=d^2z\left(\frac{1}{2}x_m\overline{}x^m+p_\alpha \overline{}\theta ^\alpha \right)$$ (3.2) where the fermionic conjugate momenta $`p_\alpha `$ are independent variables. This approach was shown to work for quantizing the superparticle, but not the superstring, since its correct physical spectrum was never obtained. We should note that the action (3.2) is quadratic and therefore it is immediate to determine the OPE’s between the free fields. Quantization of this theory is as simple as the in RNS formalism. In the next section we will see how Siegel action (3.2) is the starting point for the construction of a formalism for the superstring that allows for a covariant quantization and describes the same physics as the RNS and GS strings. #### 3.1.2 Pure spinor superstring basics From now on I will use the Weyl representation of the $`32\times 32`$ ten-dimensional Dirac matrices, where they are off-diagonal and $`\gamma _{\alpha \beta }^m`$ and $`\gamma _{\widehat{\alpha }\widehat{\beta }}^m`$ are the real symmetric $`16\times 16`$ off-diagonal blocks, satisfing the Fierz identities $`\gamma _{m\alpha (\beta }\gamma _{\gamma \delta )}^m=0`$. Useful properties to keep in mind are the following. Every symmetric bispinor can be decomposed in terms of a vector and a five form as $$f_{\alpha \beta }=\gamma _{\alpha \beta }^mf_m+\gamma _{\alpha \beta }^{mnpqr}f_{mnpqr}$$ (3.3) while every antisymmetric bispinor can be decomposed in terms of a three form as $$\stackrel{~}{f}_{\alpha \beta }=\gamma _{\alpha \beta }^{mnp}\stackrel{~}{f}_{mnp}$$ (3.4) Our conventions for $`d=10`$ $`N=2`$ superspace covariant derivatives and supersymmetry charges are $`D_\alpha =_\alpha +{\displaystyle \frac{1}{2}}(\gamma ^m\theta )_\alpha _m,Q_\alpha =_\alpha {\displaystyle \frac{1}{2}}(\gamma ^m\theta )_\alpha _m,`$ (3.6) $`\widehat{D}_{\widehat{\alpha }}=_{\widehat{\alpha }}+{\displaystyle \frac{1}{2}}(\gamma ^m\widehat{\theta })_{\widehat{\alpha }}_m,\widehat{Q}_{\widehat{\alpha }}=_{\widehat{\alpha }}{\displaystyle \frac{1}{2}}(\gamma ^m\theta )_{\widehat{\alpha }}_m,`$ which satisfy $`\{D_\alpha ,D_\beta \}=\gamma _{\alpha \beta }^m_m,\{\widehat{D}_{\widehat{\alpha }},\widehat{D}_{\widehat{\beta }}\}=\gamma _{\widehat{\alpha }\widehat{\beta }}^m_m,\{D_\alpha ,\widehat{D}_{\widehat{\beta }}\}=0`$ (3.7) $`\{D_\alpha ,Q_\beta \}=0,\{\widehat{D}_{\widehat{\alpha }},\widehat{Q}_{\widehat{\beta }}\}=0.`$ (3.8) Berkovits completed Siegel action (3.2) by adding some missing worldsheet ghost degrees of freedom. The evaluation of what’s missing in Siegel approach can be achieved by “counting”. Siegel action (3.2) gives the free-field OPE’s $`x^m(y)x^n(w)2\eta ^{mn}\mathrm{log}\|yw\|`$ (3.9) $`p_\alpha (y)\theta ^\beta (w)\delta _\alpha ^\beta (yw)^1`$ (3.10) From this, we can determine the contributions of the different fields to the conformal anomaly. Since fermionic fields contribute -32 and bosonic ones +10, the missing fields should contribute +22. Moreover, one can consider the contribution to the Lorentz current coming from fermionic degrees of freedom $`M_{mn}=\frac{1}{2}p\gamma _{mn}\theta `$ and compare to the analogous term in RNS formalism $`M_{mn}=\psi _m\psi _n`$. The current-current OPE’s are similar, except for the coefficient of the double pole term, which is +4 in Siegel case and +1 in RNS case. Therefore, the missing ghost variables should contribute to the Lorentz current in a way to produce a -3 in the double pole. Indeed, Berkovits found that an irreducible representation of $`SO(9,1)`$ with these characteristics exists. This is a bosonic pure-spinor satisfying $$\lambda \gamma ^m\lambda =0$$ (3.11) To solve this constraint and find the free ghost fields, one has to break the manifest Lorentz covariance to a $`U(5)`$ subgroup of (Wick rotated) $`SO(10)`$. In terms of this parametrization one can write down the ghost-field action, check that the OPE of the ghost contribution to the Lorentz current has a -3 coefficient in the double pole and that the stress tensor has central charge +22, as required. Apparently, one goes back to the old problem of the lack of manifest Lorentz covariance. However, one can formally write down an action in the form $$S=d^2z\left(\frac{1}{2}x_m\overline{}x^m+p_\alpha \overline{}\theta ^\alpha +w_\alpha \overline{}\lambda ^\alpha \right)+\mathrm{right}\mathrm{moving}$$ (3.12) where the independent conjugate momenta $`w_\alpha `$ of the ghost field $`\lambda ^\alpha `$ have been introduced, and then “remember” that the $`\lambda `$ fields are constrained by equation (3.11). Both the action and the pure spinor constraint are manifestly Lorentz covariant. The problem is how to deal with constrained fields in a path integral approach. When there are first class constraints in a theory, a BRST quantization procedure can be applied and the BRST charge is constructed out of the constraints themselves multiplied by ghost fields. When the constraints are second class, this does not work because the BRST charge one obtains is not nilpotent. In Berkovits approach to the superstring, the (left-moving) BRST-like charge is defined as $$Q=𝑑z\lambda ^\alpha d_\alpha $$ (3.13) where $`d_\alpha `$ is the constraint of the GS superstring that in this formalism plays the role of the supersymmetric version of the fermionic conjugate momentum $`p_\alpha `$ $$d_\alpha =p_\alpha \frac{1}{2}x^m(\gamma _m\theta )_\alpha \frac{1}{8}(\gamma ^m\theta )_\alpha (\theta \gamma _m\theta )$$ (3.14) Therefore, this construction of $`Q`$ could be reminiscent of some sort of BRST quantization of the GS superstring. From the free OPE’s (3.10) one can compute the OPE’s between the composite variables $`d_\alpha `$ and find $$d_\alpha (y)d_\beta (w)(yw)^1\gamma _{\alpha \beta }^m\mathrm{\Pi }_m(w)$$ (3.15) Therefore the BRST charge (3.13) is nilpotent because of the pure spinor condition (3.11). Also because of the pure spinor condition, one observes that the ghost conjugate momentum $`w_\alpha `$ can only appear in combinations that are invariant under the gauge transformation $$\delta w_\alpha =\mathrm{\Lambda }_m(\gamma ^m\lambda )_\alpha $$ (3.16) with arbitrary $`\mathrm{\Lambda }_m`$. These gauge-invariant combinations are the pure spinor contribution to the Lorentz current $`N_{mn}=\frac{1}{2}:(w\gamma _{mn}\lambda ):`$ and the ghost number current $`J=:w_\alpha \lambda ^\alpha :`$. In the superparticle case, the BRST charge (3.13) and pure spinor condition (3.11) can be obtained by an honest, although unusual, gauge fixing procedure starting from the Brink-Schwarz action in semi-light-cone gauge rewritten in the Siegel formalism. Unfortunately, no analogous procedure works for the superstring. Following the usual prescription of the BRST quantization rules, we could start from the GS superstring and define the quantum action as follows $$S_0=S_{GS}+Qd^2zw_\alpha \overline{}\theta ^\alpha $$ (3.17) where $`S_{GS}`$ is the Green-Schwarz action in conformal gauge . By moving on to a Siegel description for fermionic fields and by explicitly writing down all the contributions to (3.17), one obtains (3.12). Even if this looks like the usual BRST procedure, we have to notice that the BRST-like operator $`Q`$ is nilpotent up to gauge transformations (3.16). This compensates the fact that the Green-Schwarz action is not invariant under BRST transformations. In addition, we can always add BRST invariant terms to the action. However, there is no procedure to get (3.17) from an honest gauge fixing of the Green-Schwarz action (a suggestion is given in ). Now I’m going to discuss pure spinor superstring vertex operators. I will first derive the open superstring vertices for simplicity. Closed superstring vertices will be studied in much detail in section 3.1.3. Since in the open string case vertices are to be inserted on the boundary of the worldsheet, where the boundary condition $`\theta =\widehat{\theta }\|_{z=\overline{z}}`$ holds, they can be expressed in terms of the left-moving fermions only (or, more correctly, in terms of the linear combination $`\theta _+=\frac{1}{\sqrt{2}}(\theta +\widehat{\theta })`$ and the corresponding one for the fermionic momenta). Physical states for the open superstring are defined as ghost-number one states in the cohomology of $`Q`$, defined in (3.13), with $`\lambda `$ satisfying the pure spinor condition (3.11). Open superstring vertex operators with $`(\mathrm{mass})^2=\frac{n}{2}`$ are constructed out of the fields $`x^m`$, $`\theta ^\alpha `$, $`d_\alpha `$, $`\lambda ^\alpha `$ and the gauge invariant objects $`N_{mn}`$ and $`J`$ containing the ghost momenta. They are obtained as the generic combinations with ghost number one and conformal weight $`n`$ at zero momentum. Since the composite objects $`d_\alpha `$, $`N_{mn}`$ and $`J`$ carry conformal weight one and $`\lambda ^\alpha `$ carries ghost number one, it is clear that for instance the most general vertex operator at $`(\mathrm{mass})^2=0`$ is $$𝒱^{(1)}=\lambda ^\alpha A_\alpha (x,\theta )$$ (3.18) By requiring $`Q𝒱^{(1)}=0`$, one obtains the equations of motion for the spinor superfield $`A_\alpha (x,\theta )`$. By making use of the OPE $$d_\alpha (y)A_\beta (x,\theta )(w)D_\alpha A_\beta (w)$$ (3.19) one finds that the superfield $`A_\alpha (x,\theta )`$ must satisfy the equation $`\lambda ^\alpha \lambda ^\beta D_\alpha A_\beta =0`$. Because of the property (3.3) and the pure spinor condition (3.11), this is equivalent to $$\gamma _{mnpqr}^{\alpha \beta }D_\alpha A_\beta =0$$ (3.20) It can be shown that these are the superMaxwell equations written in terms of a spinor superfield. Equation $`QU=0`$ is invariant under the gauge transformation $`\delta U=Q\mathrm{\Omega }`$ where $`\mathrm{\Omega }`$ is a generic scalar superfield. Indeed, this implies the gauge transformation for the spinor superfield $`\delta A_\alpha =D_\alpha \mathrm{\Omega }`$, which is the expected gauge transformation of superMaxwell theory. Going on to the next mass level, $`(\mathrm{mass})^2=\frac{1}{2}`$, one finds that the most general vertex operator is $`𝒱_1^{(1)}`$ $`=`$ $`\lambda ^\alpha A_\alpha (x,\theta )+:\theta ^\beta \lambda ^\alpha B_{\alpha \beta }(x,\theta ):+:d_\beta \lambda ^\alpha C_\alpha ^\beta (x,\theta ):`$ (3.21) $`+`$ $`:\mathrm{\Pi }^m\lambda ^\alpha H_{m\alpha }(x,\theta ):+:J\lambda ^\alpha E_\alpha (x,\theta ):+:N^{mn}\lambda ^\alpha F_{\alpha mn}(x,\theta ):`$ (3.23) Cohomology equations and gauge transformations imply that the superfields appearing in the vertex describe a spin two multiplet. The integrated massless open superstring vertex operator $`𝑑z𝒱^{(0)}`$ can be obtained by making use of the cohomology descent equation $$[Q,𝒱^{(0)}]=𝒱^{(1)}$$ (3.24) $`𝒱^{(0)}`$ is expanded in terms of the 1-forms $`𝐗=(\theta ^\alpha ,\mathrm{\Pi }^m,d_\alpha ,\frac{1}{2}N_{mn})`$ as follows $$𝒱^{(0)}=\theta ^\alpha A_\alpha (x,\theta )+\mathrm{\Pi }^mA_m(x,\theta )+d_\alpha W^\alpha (x,\theta )+\frac{1}{2}N_{mn}F^{mn}$$ (3.25) The descent equation (3.24) is satisfied when the superfields $`A_\alpha `$, $`A_m`$, $`W^\alpha `$ and $`F_{mn}`$ are governed by the superMaxwell equations $`D_\alpha A_\beta +D_\beta A_\alpha \gamma _{\alpha \beta }^mA_m=0`$ (3.26) $`D_\alpha A_m_mA_\alpha \gamma _{m\alpha \beta }W^\beta =0`$ (3.27) $`D_\alpha W^\beta {\displaystyle \frac{1}{4}}(\gamma ^{mn})_\alpha ^\beta F_{mn}`$ (3.28) $`\lambda ^\alpha \lambda ^\beta (\gamma _{mn})_\beta ^\gamma D_\alpha F^{mn}=0`$ (3.29) The last equation is redundant, since it is implied by the previous one and by the pure spinor condition (3.11). The vertex (3.25) was first found by Siegel in , except for the pure spinor term, by making use of superspace arguments. All this discussion can be easily generalized to type II closed superstrings. The field content is $`x^m`$ where $`m=0,\mathrm{},9`$, two Majorana-Weyl spinors $`\theta ^\alpha `$, $`\widehat{\theta }^{\widehat{\alpha }}`$ with $`\alpha =\widehat{\alpha }=1,\mathrm{},16`$ (with opposite or same chirality depending whether one is in IIA or IIB case), their conjugate momenta $`p_\alpha `$, $`\widehat{p}_{\widehat{\alpha }}`$, two ghosts $`\lambda ^\alpha `$, $`\widehat{\lambda }^{\widehat{\alpha }}`$ satisfying the pure spinor conditions $$\lambda \gamma ^m\lambda =0,\widehat{\lambda }\gamma ^m\widehat{\lambda }=0,$$ and the corresponding conjugate momenta $`w_\alpha `$, $`\widehat{w}_{\widehat{\alpha }}`$. Again, supersymmetric versions of the fermionic conjugate momenta can be introduced as follows $`d_\alpha =p_\alpha {\displaystyle \frac{1}{2}}x^m(\gamma _m\theta )_\alpha {\displaystyle \frac{1}{8}}(\gamma ^m\theta )_\alpha (\theta \gamma _m\theta ),`$ (3.30) $`\widehat{d}_{\widehat{\alpha }}=\widehat{p}_{\overline{\alpha }}{\displaystyle \frac{1}{2}}\overline{}x^m(\gamma _m\widehat{\theta })_{\widehat{\alpha }}{\displaystyle \frac{1}{8}}(\gamma ^m\widehat{\theta })_{\widehat{\alpha }}(\widehat{\theta }\gamma _m\overline{}\widehat{\theta }),`$ (3.31) The BRST operators are defined by $$Q_L=𝑑z\lambda ^\alpha d_\alpha ,Q_R=𝑑\overline{z}\widehat{\lambda }^{\widehat{\alpha }}\widehat{d}_{\widehat{\alpha }}.$$ (3.32) which satisfy $$Q_L^2=𝑑z\lambda \gamma ^m\lambda \mathrm{\Pi }_m,[Q_L,Q_R]=0,Q_R^2=𝑑\overline{z}\widehat{\lambda }\gamma ^m\widehat{\lambda }\widehat{\mathrm{\Pi }}_m,$$ (3.33) where $`\mathrm{\Pi }_z^m`$ and $`\widehat{\mathrm{\Pi }}_{\overline{z}}^m`$ are the left- and right-moving supersymmetrized bosonic momenta $$\mathrm{\Pi }_z^m=x^m+\frac{1}{2}\theta \gamma ^m\theta ;\widehat{\mathrm{\Pi }}_{\overline{z}}^m=\overline{}x^m+\frac{1}{2}\widehat{\theta }\gamma ^m\overline{}\widehat{\theta }$$ (3.34) Due to pure spinor constraints, the BRST charges are nilpotent up to gauge transformations of $`w_\alpha `$, $`\widehat{w}_{\widehat{\alpha }}`$, given by $$\delta _Lw_\alpha =\mathrm{\Lambda }_m(\gamma ^m\lambda )_\alpha ,\delta _R\widehat{w}_\alpha =\widehat{\mathrm{\Lambda }}_m(\gamma ^m\widehat{\lambda })_\alpha .$$ (3.35) with arbitrary local parameters $`\mathrm{\Lambda }_m`$ and $`\widehat{\mathrm{\Lambda }}_m`$. Gauge invariant operators are \| $`J_L=:w_\alpha \lambda ^\alpha :,`$ \| $`J_R=:\widehat{w}_\alpha \widehat{\lambda }^\alpha :,`$ \| \| --- \| --- \| \| $`N_L=\frac{1}{2}:w\gamma ^{mn}\lambda :,`$ \| $`N_R=\frac{1}{2}:\widehat{w}\gamma ^{mn}\widehat{\lambda }:,`$ \| (3.36) By formally following the usual prescription of the BRST quantization rules, we can define the quantum action starting from the GS superstring in Siegel formalism as follows $$S_0=S_{GS}+Q_Ld^2zw_\alpha \overline{}\theta ^\alpha +Q_Rd^2z\widehat{w}_{\widehat{\alpha }}\widehat{\theta }^{\widehat{\alpha }}.$$ (3.37) By exploiting the different contributions in (3.37), we obtain $$S_0=d^2z\left(\frac{1}{2}x^m\overline{}x_m+p_\alpha \overline{}\theta ^\alpha +\widehat{p}_{\widehat{\alpha }}\widehat{\theta }^{\widehat{\alpha }}+w_\alpha \overline{}\lambda ^\alpha +\widehat{w}_{\widehat{\alpha }}\widehat{\lambda }^{\widehat{\alpha }}\right),$$ (3.38) which is BRST invariant and invariant under the gauge transformation (3.35) if the spinors $`\lambda ^\alpha ,\widehat{\lambda }^{\widehat{\alpha }}`$ are pure. The action is also invariant under the $`N=2`$ supersymmetry transformations generated by $`Q_ϵ=ϵ^\alpha 𝑑zq_\alpha +\widehat{ϵ}^{\widehat{\alpha }}𝑑\overline{z}\widehat{q}_{\widehat{\alpha }}`$ where the explicit expressions for the supersymmetry currents are $`q_\alpha =p_\alpha +{\displaystyle \frac{1}{2}}x^m(\gamma _m\theta )_\alpha +{\displaystyle \frac{1}{24}}(\theta \gamma ^m\theta )(\gamma _m\theta )_\alpha ,`$ (3.39) $`\widehat{q}_{\widehat{\alpha }}=\widehat{p}_{\overline{\alpha }}+{\displaystyle \frac{1}{2}}\overline{}x^m(\gamma _m\widehat{\theta })_{\widehat{\alpha }}+{\displaystyle \frac{1}{24}}(\widehat{\theta }\gamma ^m\overline{}\widehat{\theta })(\gamma _m\widehat{\theta })_{\widehat{\alpha }}.`$ (3.40) It is interesting to note that these do not anticommute with the BRST operators $`Q_L`$ and $`Q_R`$, since $$[Q_L,q_\alpha ]=\chi _\alpha ,[Q_R,\widehat{q}_{\widehat{\beta }}]=\overline{}\widehat{\chi }_{\widehat{\beta }}$$ (3.41) where $`\chi _\alpha `$ and $`\widehat{\chi }_{\widehat{\beta }}`$ are the BRST-invariant quantities $$\chi _\alpha \frac{1}{3}(\lambda \gamma ^m\theta )(\gamma _m\theta )_\alpha ,\widehat{\chi }_{\widehat{\beta }}=\frac{1}{3}(\widehat{\lambda }\gamma ^p\widehat{\theta })(\gamma _p\widehat{\theta })_{\widehat{\beta }}$$ (3.42) We also introduce the Lorentz currents $`L^{mn}={\displaystyle \frac{1}{2}}:x^{[m}x^{n]}:+{\displaystyle \frac{1}{2}}:(p\gamma ^{mn}\theta ):+:N^{mn}:,`$ (3.43) $`\widehat{L}^{pq}={\displaystyle \frac{1}{2}}:\overline{}x^{[p}x^{q]}:+{\displaystyle \frac{1}{2}}:(\widehat{p}\gamma ^{pq}\widehat{\theta }):+:\widehat{N}^{pq}:,`$ (3.44) which satisfy the following commutation relations with the BRST charges $$[Q_L,L^{mn}]=𝒢^{mn};[Q_R,\widehat{L}^{pq}]=\overline{}\widehat{𝒢}^{pq}$$ (3.45) where $$𝒢^{mn}=\frac{1}{4}(\theta \gamma ^r\lambda )(\delta _r^{[m}x^{n]}+\frac{1}{4}(\theta \gamma _r\gamma ^{mn}\theta ));\widehat{𝒢}^{pq}=\frac{1}{4}(\widehat{\theta }\gamma ^r\widehat{\lambda })(\delta _r^{[p}x^{q]}+\frac{1}{4}(\widehat{\theta }\gamma _r\gamma ^{pq}\widehat{\theta }))$$ (3.46) are BRST invariant. By using the equations of motion from (3.38) it is easy to show that the currents $`q_\alpha `$, $`\widehat{q}_{\widehat{\beta }}`$, $`\lambda ^\alpha d_\alpha `$, $`\widehat{\lambda }^{\widehat{\beta }}\widehat{d}_\beta `$, $`L^{mn}`$ and $`\widehat{L}^{pq}`$ are holomophic and anti-holomorphic, respectively. In the following section I will describe in detail type II vertex operators, their descent equations and the corresponding superfield equations of motion and gauge transformations. These closed string vertices are as usual obtained by taking the left-right product of the open superstring vertices I described in the present section. #### 3.1.3 Type II superstring vertex operators In this section I will describe in detail the construction of the closed superstring ghost number $`(1,1)`$ local vertex operator $`𝒱^{(1,1)}`$ and of the integrated vertex operators $`𝑑z𝒱_z^{(0,1)}`$, $`𝑑\overline{z}𝒱_{\overline{z}}^{(1,0)}`$, and $`𝑑zd\overline{z}𝒱_{z\overline{z}}^{(0,0)}`$, related to it by the closed string descent equations. Introducing the notation $`𝒪_{c,d}^{(a,b)}`$ for local vertex operators with ghost number $`a(b)`$ in the left (right) sector and (anti)holomorphic indices $`c(d)`$, we identify $`𝒪_{0,0}^{(1,1)}`$ $`=`$ $`𝒱^{(1,1)},`$ (3.47) $`𝒪_{1,0}^{(0,1)}=𝒱_z^{(0,1)}dz`$ , $`𝒪_{0,1}^{(1,0)}=𝒱_{\overline{z}}^{(1,0)}d\overline{z},`$ (3.48) $`𝒪_{1,1}^{(0,0)}`$ $`=`$ $`𝒱_{z\overline{z}}^{(0,0)}dzd\overline{z}.`$ (3.49) The descent equations read<sup>1</sup><sup>1</sup>1Here we use the square brackets to denote both commutation and anti-commutation relations. The difference is established by the nature of the operators involved in the relations. $$[Q_L,𝒪_{c,d}^{(a,b)}]=𝒪_{c1,d}^{(a+1,b)},[Q_R,𝒪_{c,d}^{(a,b)}]=\overline{}𝒪_{c,d1}^{(a,b+1)},$$ (3.50) where $`=dz_z`$ and $`\overline{}=d\overline{z}_{\overline{z}}`$ are the holomorphic and antiholomorphic differentials. $`Q_L`$ and $`Q_R`$ are the BRST charges for holomorphic and antiholomorphic sectors we introduced in (3.32). More explicitly, at the first level we have $$[Q_L,𝒱^{(1,1)}]=0,[Q_R,𝒱^{(1,1)}]=0,$$ (3.51) while at the next level we get $$[Q_L,𝒱_z^{(0,1)}]=_z𝒱^{(1,1)},[Q_R,𝒱_z^{(0,1)}]=0,$$ (3.52) $$[Q_R,𝒱_{\overline{z}}^{(1,0)}]=_{\overline{z}}𝒱^{(1,1)},[Q_L,𝒱_{\overline{z}}^{(1,0)}]=0,$$ (3.53) and, finally, $$[Q_L,𝒱_{z\overline{z}}^{(0,0)}]=_z𝒱_{\overline{z}}^{(1,0)},[Q_R,𝒱_{z\overline{z}}^{(0,0)}]=_{\overline{z}}𝒱_z^{(0,1)}.$$ (3.54) The vertex operators $`𝒪_{c,d}^{(a,b)}`$ are to be expanded in powers of ghost fields $`\lambda ^\alpha `$ and $`\widehat{\lambda }^{\widehat{\alpha }}`$ or in powers of the supersymmetric holomorphic and antiholomorphic 1-forms $$𝐗_z=(_z\theta ^\alpha ,\mathrm{\Pi }_z^m,d_{z\alpha },\frac{1}{2}N_z^{mn}),\widehat{𝐗}_{\overline{z}}=(_{\overline{z}}\widehat{\theta }^{\widehat{\beta }},\widehat{\mathrm{\Pi }}_{\overline{z}}^p,\widehat{d}_{\overline{z}\widehat{\beta }},\frac{1}{2}\widehat{N}_{\overline{z}}^{pq}).$$ (3.55) The explicit expressions of these 1-form operators in terms of sigma model fields are given in (3.30, 3.34, 3.36). The coefficients are superfields of the coordinates $`x^m`$, $`\theta ^\alpha `$ and $`\widehat{\theta }^{\widehat{\alpha }}`$. A further relation is obtained by acting from the left on the first equation of (3.54) with $`Q_R`$ or on the second with $`Q_L`$. Using eqs. (3.52), one obtains $$[Q_R,[Q_L,𝒱_{z\overline{z}}^{(0,0)}]]=_z_{\overline{z}}𝒱^{(1,1)},$$ (3.56) which is the closed string analogue of (3.24) and relates the integrated vertex $`𝒱_{z\overline{z}}^{(0,0)}`$ to the unintegrated one $`𝒱^{(1,1)}`$ . Equations (3.51, 3.52, 3.54) are invariant under the gauge transformations given by $$\delta 𝒱^{(1,1)}=[Q_L,\mathrm{\Lambda }^{(0,1)}]+[Q_R,\mathrm{\Lambda }^{(1,0)}]$$ (3.57) $$\delta 𝒱_{\overline{z}}^{(1,0)}=[Q_L,\tau _{\overline{z}}^{(0,0)}]+_{\overline{z}}\mathrm{\Lambda }^{(1,0)},\delta 𝒱_z^{(0,1)}=[Q_R,\tau _z^{(0,0)}]+_z\mathrm{\Lambda }^{(0,1)}$$ (3.58) $$\delta 𝒱_{z\overline{z}}^{(0,0)}=_z\tau _{\overline{z}}^{(0,0)}_{\overline{z}}\tau _z^{(0,0)}$$ (3.59) where the zero forms $`\mathrm{\Lambda }^{(0,1)}`$ and $`\mathrm{\Lambda }^{(1,0)}`$ have ghost number $`(1,0)`$ and $`(0,1)`$ and are proportional to $`\lambda ^\alpha `$ and $`\widehat{\lambda }^{\widehat{\alpha }}`$, and the coefficients are superfields. The holomorphic and antiholomorphic 1-forms $`\tau _z^{(0,0)}`$ and $`\tau _{\overline{z}}^{(0,0)}`$ are to be expanded in terms of the 1-forms $`𝐗_z`$ and $`\widehat{𝐗}_{\overline{z}}`$ given in (3.55) and coefficients are again superfields. In addition, the gauge parameters $`\mathrm{\Lambda }^{(0,1)}`$, $`\mathrm{\Lambda }^{(1,0)}`$, $`\tau _z^{(0,0)}`$ and $`\tau _{\overline{z}}^{(0,0)}`$ must satisfy the following consistency conditions $$[Q_L,\mathrm{\Lambda }^{(1,0)}]=0[Q_R,\mathrm{\Lambda }^{(0,1)}]=0,$$ (3.60) and $$[Q_L,\tau _z^{(0,0)}]+_z\mathrm{\Lambda }^{(1,0)}=0[Q_R,\tau _{\overline{z}}^{(0,0)}]+_{\overline{z}}\mathrm{\Lambda }^{(0,1)}=0.$$ (3.61) These equations resemble the descent equations for the open string vertex operator $`𝒱^{(1)}=\lambda ^\alpha A_\alpha `$, but in that case there are boundary conditions for the fermionic fields: $`\theta ^\alpha (z)=\widehat{\theta }^{\widehat{\alpha }}(z)`$ at $`z=\overline{z}`$. Equations (3.60) and (3.61) are further invariant under the gauge transformations $$\delta \mathrm{\Lambda }^{(1,0)}=[Q_L,\mathrm{{\rm Y}}^{(0,0)}],\delta \mathrm{\Lambda }^{(0,1)}=[Q_R,\widehat{\mathrm{{\rm Y}}}^{(0,0)}],$$ (3.62) $$\delta \tau _z^{(0,0)}=_z\mathrm{{\rm Y}}^{(0,0)},\delta \tau _{\overline{z}}^{(0,0)}=_{\overline{z}}\widehat{\mathrm{{\rm Y}}}^{(0,0)}.$$ (3.63) where $`\mathrm{{\rm Y}}^{(0,0)}`$ and $`\widehat{\mathrm{{\rm Y}}}^{(0,0)}`$ are generic superfields. However, consistency with (3.58) imposes $`\mathrm{{\rm Y}}^{(0,0)}=\widehat{\mathrm{{\rm Y}}}^{(0,0)}`$. The superfield $`\mathrm{{\rm Y}}^{(0,0)}`$ will be useful to define a suitable gauge fixing procedure and to take into account the reducible gauge symmetry of the NS-NS two form of 10-dimensional supergravity. To derive equations (3.52) we can view the vertex operators $`𝒱_z^{(0,1)}`$ and $`𝒱_{\overline{z}}^{(1,0)}`$ as deformations of the BRST charges $$Q_LQ_L+𝑑\overline{z}𝒱_{\overline{z}}^{(1,0)},Q_RQ_R+𝑑z𝒱_z^{(0,1)},$$ (3.64) and the vertex operator $`𝒱_{z\overline{z}}^{(0,0)}`$ as the deformation of the action $$SS+𝑑z𝑑\overline{z}𝒱_{z\overline{z}}^{(0,0)}.$$ (3.65) Eqs. (3.51) are derived by requiring the nilpotency of the new charges and the vanishing of their anticommutation relation. #### 3.1.4 Amplitudes Since the worldsheet ghost variables $`\lambda ^\alpha `$ are constrained by the pure spinor condition (3.11), it is not obvious how the define a path integral in this variables and therefore how to compute superstring amplitudes. For this reason in a different formulation was proposed where the pure spinor constraint is relaxed by adding more fields to the theory. Clearly this should be done without modifying the BRST cohomology, that was shown to reproduce the correct superstring physical spectrum . However, Berkovits recently showed that multiloop superstring amplitudes can be computed in the pure spinor formalism, by introducing an analogue of the RNS “picture changing” operators. When only tree-level amplitudes are under concern, a prescription can be given to compute them relying on properties of BRST cohomology, as shown in . The prescription given there was shown to coincide with the standard RNS one in . In terms of the vertex operators $`𝒪_{c,d}^{(a,b)}`$, the amplitudes on the sphere are defined as $$𝒜_{n+3}=𝒱^{(1,1)}(z_1,\overline{z}_1)𝒱^{(1,1)}(z_2,\overline{z}_2)𝒱^{(1,1)}(z_3,\overline{z}_3)\underset{n}{}𝑑z𝑑\overline{z}𝒱^{(0,0)}$$ (3.66) where the three unintegrated vertex operators are needed to fix the $`SL(2,𝐂)`$ invariance on the sphere. An unintegrated vertex $`𝒱^{(1,1)}(z_1,\overline{z}_1)`$ can be replaced by a product of $`(1,0)`$ and $`(0,1)`$ vertices $`𝑑z𝒱_z^{(0,1)}𝑑\overline{z}𝒱_{\overline{z}}^{(1,0)}`$ which has the same total ghost number and the same total conformal spin as the original vertex $`𝒱^{(1,1)}`$. In supersymmetry and gauge invariance were proven under the assumption that the prescription for the zero modes is the following $$𝒱^{(3,3)}=1$$ (3.67) where $$𝒱^{(3,3)}=(\lambda _0\gamma ^m\theta _0\lambda _0\gamma ^n\theta _0\lambda _0\gamma ^p\theta _0\theta _0\gamma _{mnp}\theta _0)(\widehat{\lambda }_0\gamma ^m\widehat{\theta }_0\widehat{\lambda }_0\gamma ^n\widehat{\theta }_0\widehat{\lambda }_0\gamma ^p\widehat{\theta }_0\widehat{\theta }_0\gamma _{mnp}\widehat{\theta }_0).$$ (3.68) As anticipated, this zero-mode prescription is justified by cohomological arguments. In fact, by analogy with the RNS case, one deduces that the the expectation value for the +3 ghost-number vertex operator for the Yang-Mills antighost must be fixed to one and there is a unique state of such ghost number in the pure-spinor BRST cohomology, given by $`𝒱^{(3,3)}`$. In Berkovits has given a general prescription for the computation of multiloop amplitudes in the pure spinor formalism for the superstring. Since in my work I didn’t compute amplitudes, I’m not going to give technical details about this. However, I will discuss the various difficulties that were overcome in . As outlined in section 3.1.1, the ghost variables can appear only as $`\lambda ^\alpha `$ (with conformal weight 0) or in the two gauge-invariant combinations $`N_{mn}`$ and $`J`$ (with conformal weight 1). Because of the pure spinor constraint (3.11), $`\lambda ^\alpha `$ has only eleven free components. As a result, on a genus $`g`$ surface $`\lambda ^\alpha `$ has eleven independent zero-modes and $`N_{mn}`$ and $`J`$ have $`11g`$ ones. It is not obvious how to determine a Lorentz covariant prescription to integrate over these ghost zero-modes. A Lorentz-invariant measure for $`\lambda `$ $`[𝒟\lambda ]`$ was constructed in . Moreover, it has been noted that zero modes of $`N`$ and $`J`$ are related by a constraint following from the pure spinor relation (3.11), such that all these zero-modes can be expressed in terms of ten free $`N`$ zero-modes. The measure for these free modes, $`[𝒟N]`$, is also given in . On a genus $`g`$ surface, one integrates out all the non-ghost fields and the ghost non-zero-mode fields by making use of the given OPE’s. One is then left with an expression like $$𝒜=f(\lambda ,N_1,J_1,\mathrm{},N_g,J_g)$$ (3.69) that only depends on the ghost zero-modes. Then one uses the Lorentz invariant measures to define the integration over these zero-modes $$𝒜=[𝒟\lambda ][𝒟N_1]\mathrm{}[𝒟N_g]f(\lambda ,N_1,J_1,\mathrm{},N_g,J_g)$$ (3.70) To compute the integral, as in RNS formalism , it is necessary to insert picture changing operators involving the delta-functions $`\delta (C_\alpha \lambda ^\alpha )`$, $`\delta (B_{mn}N^{mn})`$ and $`\delta (J)`$, where the constant spinor and antisymmetric tensor $`C_\alpha `$ and $`B_{mn}`$ should not affect the amplitudes (that otherwise could not be Lorentz invariant!). These picture changing operators will be called $`Y_C`$, $`Z_B`$ and $`Z_J`$ respectively. Their insertion in loop amplitudes is necessary to absorb the ghost zero-modes. Exactly as in the RNS formulation, picture changing operators must be BRST-invariant with a BRST-trivial worldsheet derivative, this second requirement needed for the amplitudes to be independent of PCO’s positions on the worldsheet. In operators with these properties were constructed and it was also shown that even if they are not supersymmetric, their variation is BRST-trivial. As a result, supersymmetry is preserved (up to surface terms). In it was also shown that the computation of tree amplitudes with the Lorentz covariant measure for pure spinor zero-modes agrees with the previous prescription obtained by cohomology arguments. To compute loop amplitudes a last ingredient is missing. In fact, in RNS formalism, in the computation of a $`g`$-loop amplitude the insertion of $`(3g3)`$ $`b`$-ghosts is necessary. The $`b`$-ghost has $`1`$ ghost number and satisfies the relation $`\{Q,b(u)\}=T(u)`$ where $`T`$ is the stress tensor. The problem is that in the pure spinor formalism the ghost conjugate momentum $`w_\alpha `$ only appears in the gauge invariant combinations $`N_{mn}`$ and $`J`$ that have both ghost-number zero. So apparently there is no candidate for an analogue of the $`b`$ ghost in this formalism. However, in it was shown that the picture raising operator $`Z_B`$, that will be present in the expression of a general multiloop amplitude, can be used to construct a nonlocal operator $`\stackrel{~}{b}_B`$, carrying ghost-number zero, such that $`\{Q,\stackrel{~}{b}_B(u,z)\}=T(u)Z_B(z)`$. From $`\stackrel{~}{b}_B(u,z)`$ it is possible to define a local $`b_B(u)`$, however for the computation of the amplitudes it will be sufficient to know the nonlocal operator (about this topic, see also ). With all these ingredients one can give the following super-Poincaré covariant prescription for the computation of a generic $`N`$-point $`g`$-loop closed string amplitude $`𝒜`$ $`={\displaystyle }d^2\tau _1\mathrm{}d^2\tau _{(3g3)}\|{\displaystyle \underset{P=1}{\overset{3g3}{}}}{\displaystyle }d^2u_P\mu _P(u_P)\stackrel{~}{b}_{B_P}(u_P,z_P)`$ (3.73) $`{\displaystyle \underset{P=3g2}{\overset{10g}{}}}Z_{B_P}(z_P){\displaystyle \underset{R=1}{\overset{g}{}}}Z_J(v_R){\displaystyle \underset{I=1}{\overset{11}{}}}Y_{C_I}(y_I)\|^2{\displaystyle \underset{T=1}{\overset{N}{}}}{\displaystyle }d^2t_T𝒱_T(t_T)`$ where $`\|\|^2`$ means left-right product, $`\tau _P`$ are the Teichmuller parameters associated to the Beltrami differentials $`\mu _P(u_P)`$ and $`𝒱_T(t_T)`$ are the dimension $`(1,1)`$ closed superstring vertex operators for the $`N`$ external states. When $`g=1`$, the general prescription (3.73) must be modified, as usual, by changing one integrated vertex operator with an unintegrated one. Even though this formalism is quite cumbersome for computing superstring amplitudes, it makes it easy to prove some general vanishing theorems, as it is somehow expected because of the great amount of manifest symmetries. Actually, the general vanishing of certain amplitudes can be deduced by just counting zero-modes. For instance, S-duality of type IIB superstrings imply that $`R^4`$ terms in the low energy effective action do not receive perturbative corrections above one loop . In RNS formalism, this was checked only up to two loops in . In the pure spinor formalism for the superstring, instead, this can be easily proven for general $`g`$ . Furthermore, it is well-known that massless N-point superstring $`g`$-loop amplitudes vanish for $`N<4`$. This is equivalent to perturbative finiteness of the theory, when unphysical divergences are not present in the interior of moduli space . This was proven in by an argument that made use of both GS and RNS formalisms. In the pure spinor formalism this can also be proven by counting zero-modes . #### 3.1.5 What we can (cannot) do with this formalism, up to now In this section I would like to summarize the results obtained with the pure spinor formalism for the superstring. First of all, I should say that in it was proven that the pure spinor BRST cohomology reproduces the correct superstring spectrum. Therefore, the pure spinor formalism describes the same physics as the RNS and the GS formalisms. In the following I will describe how the pure spinor superstring proved to be superior to analyze many aspects of string theory. As discussed in the previous sections, the main motivation to introduce a superPoincaré covariant formulation for the superstring is the number of difficulties one encounters when trying to compute a general superstring amplitude in the RNS and GS formalisms. As we have seen in the previous section, the pure spinor formalism for the superstring in principle allows to compute arbitrary $`N`$-point multiloop amplitudes (even though, in practice, only tree-level and four-point one-loop amplitudes (that can be computed also in RNS or GS formalisms) have been explicitly evaluated up to now). Furthermore, we have already stressed that a big advantage of the pure spinor formulation is the possibility to prove vanishing theorems at arbitrary order in the perturbative expansion. When I wrote my paper , in collaboration with P.A. Grassi, the prescription for computing a general multiloop amplitude was not known. However it was already clear that the computation of the amplitudes in the pure spinor formalism is rather involved, even in the tree-level case. One of the biggest problems is the complexity of the expression for unintegrated and integrated vertex operators, due to the manifest symmetries that render the formulation redundant. As we have seen in section 3.1.3, vertices are written in terms of superfields satisfying a set of linearized equations of motion, where the physical fields, such as the graviton, dilaton, R-R field strength and so on, do not appear explicitly. To write down the component expansion of the vertex, one has to solve the equations of motion, after having chosen a specific gauge. This procedure is quite complicated and determining the vertex for a given physical field is not an easy task. In , P.A. Grassi and I described an iterative procedure to compute type II vertices that eliminates auxiliary fields from the vertices and allows to determine the whole $`\theta `$ and $`\widehat{\theta }`$ vertex expansion given the physical fields. Another serious problem of RNS formalism is the difficulty in dealing with general R-R background. We have seen that vertices for closed superstring can be written and that the equations of motion and gauge trasformations for the superfields appearing in the vertices are the linearized supergravity equations and gauge transformations. Moreover, the pure spinor superstring can be naturally coupled to a general supergravity background. It has been shown in that nilpotency and holomorphicity of the pure spinor BRST charge imply the on-shell superspace constraints of the supergravity background. Aspects of the superstring in specific R-R backgrounds, such as $`AdS_5\times S^5`$ and the pp-wave were considered in . Furthermore, as we have seen in section 1.2.4, non(anti)commutative superspaces were shown to emerge when open superstrings in the presence of D-branes and R-R backgrounds are considered. The natural setting for this discussion was the ten-dimensional pure spinor superstring and its compactification on a CY three-fold . In the open string case, requiring BRST invariance also implies the correct equations of motion for the background fields. Indeed, in it was shown that classical BRST invariance of the open pure spinor superstring implies the supersymmetric Born-Infeld equations, that were first determined by making use of superembedding techniques in (for the application of the superembedding formalism to determine higher-order corrections to the effective dynamics of string/M theory branes, see also ). $`\kappa `$-symmetry in the GS string and superembedding formalism is replaced by BRST symmetry in the pure spinor formalism.To obtain the supersymmetric Born-Infeld equations from the pure spinor formalism for the superstring, one requires that the left and right pure spinor BRST currents are equal on the worldsheet boundary in the presence of the background. I would also like to say that the pure spinor formalism has been successfully used to quantize the $`d=10`$ superparticle . Moreover, by replacing ten-dimensional pure spinors with eleven-dimensional pure spinors, the formalism has been extended to the $`d=11`$ superparticle and supermembrane . The covariant prescription to compute loop amplitudes in the covariant formalism for the superstring I briefly discussed in section 3.1.4 has been generalized to the eleven-dimensional superparticle in . The supermembrane is a quite problematic object in string theory, since the impossibility of performing a covariant quantization has made it difficult to study its properties. However, the supermembrane is worth studying, since it is expected to be related to M-theory, which is the underlying eleven dimensional theory from which the nonpertubative symmetries of string theory are believed to come. In Berkovits “replaced” $`\kappa `$-symmetry with BRST symmetry as he did for the superstring. However, not all the problems of the supermembrane are solved by moving on to a pure spinor description, since the pure spinor supermembrane action is not quadratic in a flat background. In a conjecture is made that supermembrane amplitudes can however be computed and that they are M-theory scattering amplitudes which, after a suitable compactification that reduces to ten dimensions, reproduce Type IIA superstring scattering amplitudes. These amplitudes would contain non-perturbative information about the Type IIA superstring which might be useful for studying M-theory. Finally, in my paper , P. A. Grassi and I discussed an application of pure spinor techniques to the construction of off-shell vertex operators in the asymmetric picture, that could be useful to study the coupling of R-R potentials to D-branes. We also proposed an application to closed string field theory, by studying antifields in this formalism and constructing a kinetic term for the closed string field theory action that seems to respect the correct symmetries and generate the right equations. To conclude, I would like to stress that, even if the construction of the pure spinor formalism might seem “ad hoc” and not justified by an underlying general principle, it is the first successful covariant method to quantize the superstring. It has already proven to be useful to deal with many aspect of string theory that were unreachable by the other formalisms. It might be that in the end a more natural and elegant construction of the covariant superstring will be found. However, it is now very clear that the pure spinor approach at least goes in the right direction and is very effective and useful in string theory. ### 3.2 An iterative procedure to compute the vertex operators In this section, I will present the general procedure to compute pure spinor closed string vertex operators I introduced in , in collaboration with P.A. Grassi. #### 3.2.1 Warm up: The open superstring case Motivated by the increasing interest in the covariant techniques for computation of the amplitudes in string theory, in P.A. Grassi and I provided a calculation scheme for superstring vertex operators in pure spinor approach . Since the amount of symmetries that are manifest in the covariant formulation increases, also the number of auxiliary fields increases and a useful technique to compute the basic ingredients is needed. In we provided such a procedure and some applications (that will be discussed in section 3.3). First of all I will briefly review the open superstring case, to explain the main idea that will be applied in the next sections to the more complicated closed string case. In the case of the open superstring, we have seen in section 3.1.2 that the massless sector is described by a vertex operator $`𝒱^{(1)}=\lambda ^\alpha A_\alpha `$ at ghost number one, where $`\lambda ^\alpha `$ is a pure spinor satisfying (3.11) and $`A_\alpha (x,\theta )`$ is the spinorial component of the superconnection. The superfield $`A_\alpha `$ can be completely expressed in terms of the gauge field $`a_m(x)`$ and the gluino $`\psi ^\alpha (x)`$, for example as $$A_\alpha (x,\theta )=\frac{1}{2}(\gamma ^m\theta )_\alpha a_m(x)+\frac{1}{3}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma \psi ^\gamma (x)+𝒪(\theta ^3).$$ (3.74) The vertex operator $`𝒱^{(1)}`$ belongs to the cohomology of the BRST charge $`Q=𝑑\sigma \lambda ^\alpha d_{\sigma \alpha }`$, where $`d_{\sigma \alpha }`$ is defined in section 3.1.2, if and only if the components of $`A_\alpha `$ satisfy the linear Maxwell and Dirac equations $$^m(_ma_n_na_m)=0,\gamma _{\alpha \beta }^m_m\psi ^\beta =0.$$ (3.75) The contributions $`𝒪(\theta ^3)`$ are given in terms of the derivatives of $`a_m`$ and $`\psi `$ and are completely fixed by the equations of motion (3.29) given in , where $`A_m`$ is the vectorial part of the superconnection and $`D_\alpha =_\alpha +\frac{1}{2}(\gamma ^m\theta )_\alpha _m`$ is the superderivative. The lowest components of $`A_\alpha `$ in (3.74) are eliminated by a gauge fixing condition. Even though the computation of all terms in the expansion of $`A_\alpha `$ seems a straightforward procedure, technically it is rather involved. However, there exists a powerful technique which simplifies the task. The main idea is to choose a suitable gauge fixing such as for instance $$\theta ^\alpha A_\alpha (x,\theta )=0,$$ (3.76) which reduces the independent components in the superfield $`A_\alpha `$. This choice<sup>2</sup><sup>2</sup>2The following gauge condition has a counterpart in bosonic string theory: $`x^mA_m(x)=0`$. This fixes the gauge invariance under $`\delta A_m=_m\omega (x)`$ and it coincides with the Lorentz gauge in momentum space $`_{p_\mu }\stackrel{~}{A}_m=0`$. The gauge fixing yields the equation $`(1+x^n_n)A_m=x^nF_{mn}`$ which can be solved directly by inverting $`(1+x^n_n)`$ and obtaining $`A_m=^xd^{26}y[(1+y^p_p)^1(y^nF_{mn}(y)]`$. fixes part of the super-gauge transformation $`\delta A_\alpha =D_\alpha \mathrm{\Omega }`$, where $`\mathrm{\Omega }`$ is a scalar superfield with ghost number zero. To reach the gauge (3.76), we have to impose $`\theta ^\alpha (A_\alpha +\delta A_\alpha )=0`$, which implies that $`\theta ^\alpha D_\alpha \mathrm{\Omega }=\theta ^\alpha A_\alpha `$. Expanding $`\mathrm{\Omega }`$ as $`\mathrm{\Omega }=_{n0}\mathrm{\Omega }_{[\alpha _1\mathrm{}\alpha _n]}\theta ^{\alpha _1}\mathrm{}\theta ^{\alpha _n}`$, all components with $`n1`$ are fixed except the lowest component $`\mathrm{\Omega }_0`$, which corresponds to the usual bosonic gauge transformation of Maxwell theory. Acting with $`D_\alpha `$ on (3.76) and using the equations of motion (3.29), one gets the recursive relations $`(1+𝐃)A_\alpha =(\gamma ^m\theta )_\alpha A_m`$ (3.77) $`𝐃A_m=(\gamma _m\theta )_\gamma W^\gamma `$ (3.78) $`𝐃W^\alpha ={\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha F_{mn}`$ (3.79) $`𝐃F_{mn}=(\gamma _{[m}\theta )_\gamma _{n]}W^\gamma `$ (3.80) where $`𝐃\theta ^\alpha _\alpha `$. So, given the zero-order component of $`A_m`$, we can compute the order-$`\theta `$ component of $`A_\alpha `$. The same can be done for $`A_m`$, the spinorial field strength $`W^\alpha `$ and the bosonic curvature $`F_{mn}=_{[m}A_{n]}`$ making use of the other three equations. This renders the task of computing all components of $`A_\alpha `$ in terms of initial data $`A_m(x)=a_m(x)+𝒪(\theta )`$ and $`W^\alpha (x)=\psi ^\alpha (x)+𝒪(\theta )`$ a purely algebraic problem ( and ). Moreover, one is able to compute all components of the superfields appearing in the (descent) ghost-number-zero vertex operator $`𝒱_\sigma ^{(0)}`$ $$𝒱_\sigma ^{(0)}=_\sigma \theta ^\alpha A_\alpha +\mathrm{\Pi }_\sigma ^mA_m+d_{\sigma \alpha }W^\alpha +\frac{1}{2}N_\sigma ^{mn}F_{mn},$$ (3.81) which satisfies the descent equation $`[Q,𝒱_\sigma ^{(0)}]=_\sigma 𝒱^{(1)}`$. Here $`\sigma `$ is the boundary worldsheet coordinate and $`N_\sigma ^{mn}=\frac{1}{2}w_\sigma \gamma ^{mn}\lambda `$ is the pure spinor part of the Lorentz current. As we have seen in section 3.1.2, the operators $`\mathrm{\Pi }_\sigma ^m`$ and $`d_{\sigma \alpha }`$ are the supersymmetric line element and the fermionic constraint of the Green-Schwarz superstring , respectively. In my paper , we applied the same technique to IIA/IIB supergravity. Starting from the vertex operators for closed superstrings, we derived the complete set of equations from the BRST cohomology and we defined all curvatures and gauge transformations. Then, we imposed a set of gauge fixing conditions to remove the lowest components of the superfields and we derived an iterative procedure to compute all components. We showed that a further gauge fixing is needed to fix the reducible gauge symmetries and we showed that all chosen gauges can indeed be reached. The procedure for closed strings is original by itself, but, more importantly, our analysis leads to a generalization of (3.76) to all vertex operators, associated to both massless and massive states. Indeed, in we showed that the gauge fixing (3.76) can be written in terms of a new nilpotent charge $`𝒦`$ (with negative ghost number) as follows $$\{𝒦,𝒱^{(1)}\}=0.$$ (3.82) This imitates the Siegel gauge in string field theory. When restricted to massless states, this generalized gauge fixing condition reduces to the gauge fixing (3.76) for open superstrings and to the corresponding gauge fixing for closed strings. When applied to massive states, (3.82) also leads to a suitable gauge fixing. In our paper, we explicitly derived the gauge conditions for the first massive state for the open superstring. Again, (3.82) fixes all auxiliary fields in terms of the physical on-shell data and eliminates the lowest components. I would like to stress that in we only considered deformations (vertex operators) at first order in the coupling constant, neglecting the backreaction of background fields. #### 3.2.2 Linearized IIA/IIB supergravity equations In the present section we derive the equations of motion for the massless background fields in superspace from the BRST cohomology of the superstring. Let us start from the simplest equations (3.51) for the vertex $`𝒱^{(1,1)}`$ whose general expression is $$𝒱^{(1,1)}=\lambda ^\alpha A_{\alpha \widehat{\beta }}\widehat{\lambda }^{\widehat{\beta }}.$$ (3.83) The superfield $`A_{\alpha \widehat{\beta }}(x,\theta ,\widehat{\theta })`$ satisfies the equations of motion $$\gamma _{mnopq}^{\alpha \beta }D_\alpha A_{\beta \widehat{\beta }}=0,\gamma _{mnopq}^{\widehat{\alpha }\widehat{\beta }}D_{\widehat{\alpha }}A_{\alpha \widehat{\beta }}=0,$$ (3.84) where $`\gamma _{mnopq}^{\alpha \beta }`$ is the antisymmetrized product of five gamma matrices. The pure spinor conditions imply that only the 5-form parts of the $`D_\alpha A_{\beta \widehat{\beta }}`$ and $`D_{\widehat{\alpha }}A_{\alpha \widehat{\beta }}`$ are indeed constrained . By using Bianchi identities, one can show that they yield the type IIA/IIB supergravity equations of motion at the linearized level. All auxiliary fields present in the superfield $`A_{\alpha \widehat{\beta }}`$ are fixed by eqs. (3.84). As outlined before, one can use different types of vertices to simplify the computations. Integrated vertices are written in terms of a huge number of different superfields, whose components are completely fixed by the equations of motion. As a result, these vertices are quite complicated espressions. The set of superfields needed to compute $`𝒱^{(0,0)},\mathrm{},𝒱^{(1,1)}`$ can be grouped into the following matrix $$𝐀=\left[\begin{array}{cccc}A_{\alpha \widehat{\beta }}& A_{\alpha p}& E_\alpha ^{\widehat{\beta }}& \mathrm{\Omega }_{\alpha ,pq}\\ A_{m\widehat{\beta }}& A_{mp}& E_m^{\widehat{\beta }}& \mathrm{\Omega }_{m,pq}\\ E_{\widehat{\beta }}^\alpha & E_p^\alpha & P^{\alpha \widehat{\beta }}& C_{pq}^\alpha \\ \mathrm{\Omega }_{mn,\widehat{\beta }}& \mathrm{\Omega }_{mn,p}& C_{mn}^{\widehat{\beta }}& S_{mn,pq}\end{array}\right]$$ (3.85) The first components of $`A_{mp}`$, $`E_m^{\widehat{\beta }}`$, $`E_p^\alpha `$ and $`P^{\alpha \widehat{\beta }}`$ are identified with the supergravity fields as follows $$A_{mp}=g_{mp}+b_{mp}+\eta _{mp}\varphi +𝒪(\theta ,\widehat{\theta }),$$ (3.86) $$E_m^{\widehat{\beta }}=\psi _m^{\widehat{\beta }}+𝒪(\theta ,\widehat{\theta }),E_p^\alpha =\psi _p^\alpha +𝒪(\theta ,\widehat{\theta }),$$ (3.87) $$P^{\alpha \widehat{\beta }}=f^{\alpha \widehat{\beta }}+𝒪(\theta ,\widehat{\theta }).$$ (3.88) The fields $`g_{mn}`$, $`b_{mn}`$, $`\varphi `$, $`\psi _p^\alpha `$, $`\psi _m^{\widehat{\beta }}`$ and $`f^{\alpha \widehat{\beta }}`$ are the graviton, the NS-NS two-form, the dilaton, the two gravitinos (the gamma-traceless part of $`\psi _p^\alpha `$, $`\psi _m^{\widehat{\beta }}`$), the two dilatinos (the gamma-trace part of $`\psi _p^\alpha `$, $`\psi _m^{\widehat{\beta }}`$) and the RR field strengths. IIA and IIB differ in the chirality of the two spinorial indices $`\alpha `$ and $`\widehat{\alpha }`$. This changes the type of RR fields present in the spectrum. The first components of the superfields $`\mathrm{\Omega }_{m,pq}`$ ($`\mathrm{\Omega }_{mn,p}`$), $`C_{mn}^\beta `$ ($`C_{pq}^\alpha `$) and $`S_{mn,pq}`$ are identified with the linearized gravitational connection $`\mathrm{\Gamma }_{rs}^t`$, the curvature of the gravitinos and the linearized Riemann tensor, respectively. The remaining superfields are the spinorial partners of the above superfields. Those constraints are given in terms of the spinorial components $`A_{\alpha \widehat{\beta }}`$, $`A_{\alpha p}`$, $`E_\alpha ^{\widehat{\beta }}`$ and $`\mathrm{\Omega }_{\alpha ,pq}`$. The structure of superspace formulation of type IIA and IIB supergravity in the present framework is also discussed in . Given the vectors $`𝐗_z`$ and $`\widehat{𝐗}_{\overline{z}}`$ (see (3.55)) we can explicitly write the vertex operator $`𝒱_{z\overline{z}}^{(0,0)}=𝐗_z^T𝐀\widehat{𝐗}_{\overline{z}}`$ as $`𝒱_{z\overline{z}}^{(0,0)}`$ $`=`$ $`_z\theta ^\alpha A_{\alpha \widehat{\beta }}_{\overline{z}}\widehat{\theta }^{\widehat{\beta }}+_z\theta ^\alpha A_{\alpha p}\widehat{\mathrm{\Pi }}_{\overline{z}}^p+\mathrm{\Pi }_z^mA_{m\widehat{\beta }}_{\overline{z}}\widehat{\theta }^{\widehat{\beta }}+\mathrm{\Pi }_z^mA_{mp}\widehat{\mathrm{\Pi }}_{\overline{z}}^p`$ (3.89) $`+`$ $`d_{z\alpha }E_{\widehat{\beta }}^\alpha _{\overline{z}}\widehat{\theta }^{\widehat{\beta }}+d_{z\alpha }E_p^\alpha \widehat{\mathrm{\Pi }}_{\overline{z}}^p+_z\theta ^\alpha E_\alpha ^{\widehat{\beta }}\widehat{d}_{\overline{z}\widehat{\beta }}+\mathrm{\Pi }_z^mE_m^{\widehat{\beta }}\widehat{d}_{\overline{z}\widehat{\beta }}`$ (3.90) $`+`$ $`d_{z\alpha }P^{\alpha \widehat{\beta }}\widehat{d}_{\overline{z}\widehat{\beta }}+{\displaystyle \frac{1}{2}}N_z^{mn}\mathrm{\Omega }_{mn,\widehat{\beta }}_{\overline{z}}\widehat{\theta }^{\widehat{\beta }}+{\displaystyle \frac{1}{2}}N_z^{mn}\mathrm{\Omega }_{mn,p}\widehat{\mathrm{\Pi }}_{\overline{z}}^p`$ (3.91) $`+`$ $`{\displaystyle \frac{1}{2}}_z\theta ^\alpha \mathrm{\Omega }_{\alpha ,pq}\widehat{N}_{\overline{z}}^{pq}+{\displaystyle \frac{1}{2}}\mathrm{\Pi }_z^m\mathrm{\Omega }_{m,pq}\widehat{N}_{\overline{z}}^{pq}+{\displaystyle \frac{1}{2}}N_z^{mn}C_{mn}^{\widehat{\beta }}\widehat{d}_{\overline{z}\widehat{\beta }}`$ (3.92) $`+`$ $`{\displaystyle \frac{1}{2}}d_{z\alpha }C_{pq}^\alpha \widehat{N}_{\overline{z}}^{pq}+{\displaystyle \frac{1}{4}}N_z^{mn}S_{mn,pq}\widehat{N}_{\overline{z}}^{pq}`$ (3.93) From equations (3.51), (3.52), (3.54) and (3.56) in the previous section we derive the complete set of equations for the background fields \| $`(\frac{1}{2},\frac{1}{2},\frac{1}{2})`$ \| $`D_\alpha A_{\beta \widehat{\gamma }}+D_\beta A_{\alpha \widehat{\gamma }}\gamma _{\alpha \beta }^mA_{m\widehat{\gamma }}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}A_{\beta \widehat{\gamma }}+\widehat{D}_{\widehat{\gamma }}A_{\beta \widehat{\alpha }}\gamma _{\widehat{\alpha }\widehat{\gamma }}^mA_{\beta m}=0`$ \| \| --- \| --- \| --- \| \| $`(\frac{1}{2},\frac{1}{2},1)`$ \| $`D_\alpha A_{m\widehat{\beta }}_mA_{\alpha \widehat{\beta }}\gamma _{m\alpha \gamma }E_{\widehat{\beta }}^\gamma =0`$ \| $`\widehat{D}_{\widehat{\alpha }}A_{\beta p}_pA_{\beta \widehat{\alpha }}\gamma _{p\widehat{\alpha }\widehat{\gamma }}E_\beta ^{\widehat{\gamma }}=0`$ \| \| $`(\frac{1}{2},\frac{1}{2},1)`$ \| $`D_\alpha A_{\beta p}+D_\beta A_{\alpha p}\gamma _{\alpha \beta }^mA_{mp}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}A_{m\widehat{\beta }}+\widehat{D}_{\widehat{\beta }}A_{m\widehat{\alpha }}+\gamma _{\widehat{\alpha }\widehat{\beta }}^pA_{mp}=0`$ \| \| $`(\frac{1}{2},\frac{1}{2},\frac{3}{2})`$ \| $`D_\alpha E_{\widehat{\gamma }}^\beta \frac{1}{4}\left(\gamma ^{mn}\right)_\alpha ^\beta \mathrm{\Omega }_{mn,\widehat{\gamma }}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}E_\beta ^{\widehat{\gamma }}\frac{1}{4}\left(\gamma ^{pq}\right)_{\widehat{\alpha }}^{\widehat{\gamma }}\mathrm{\Omega }_{\beta ,pq}=0`$ \| \| $`(\frac{1}{2},\frac{1}{2},\frac{3}{2})`$ \| $`D_\alpha E_\beta ^{\widehat{\gamma }}+D_\beta E_\alpha ^{\widehat{\gamma }}\gamma _{\alpha \beta }^mE_m^{\widehat{\gamma }}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}E_{\widehat{\gamma }}^\beta +\widehat{D}_{\widehat{\gamma }}E_{\widehat{\alpha }}^\beta \gamma _{\widehat{\alpha }\widehat{\gamma }}^pE_p^\beta =0`$ \| \| $`(\frac{1}{2},1,1)`$ \| $`D_\alpha A_{mp}_mA_{\alpha p}\gamma _{m\alpha \gamma }E_p^\gamma =0`$ \| $`\widehat{D}_{\widehat{\alpha }}A_{mp}+_pA_{m\widehat{\alpha }}+\gamma _{p\widehat{\alpha }\widehat{\beta }}E_m^{\widehat{\beta }}=0`$ \| \| $`(\frac{1}{2},\frac{3}{2},1)`$ \| $`D_\alpha E_p^\beta \frac{1}{4}\left(\gamma ^{mn}\right)_\alpha ^\beta \mathrm{\Omega }_{mn,p}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}E_m^{\widehat{\beta }}+\frac{1}{4}\mathrm{\Omega }_{m,pq}\left(\gamma ^{pq}\right)_{\widehat{\alpha }}^{\widehat{\beta }}=0`$ \| \| $`(\frac{1}{2},\frac{3}{2},1)`$ \| $`D_\alpha E_m^{\widehat{\beta }}_mE_\alpha ^{\widehat{\beta }}\gamma _{m\alpha \gamma }P^{\gamma \widehat{\beta }}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}E_p^\beta _pE_{\widehat{\alpha }}^\beta \gamma _{p\widehat{\alpha }\widehat{\gamma }}P^{\beta \widehat{\gamma }}=0`$ \| \| $`(\frac{1}{2},\frac{1}{2},2)`$ \| $`D_\alpha \mathrm{\Omega }_{\beta ,pq}+D_\beta \mathrm{\Omega }_{\alpha ,pq}\gamma _{\alpha \beta }^m\mathrm{\Omega }_{m,pq}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}\mathrm{\Omega }_{mn,\widehat{\beta }}+\widehat{D}_{\widehat{\beta }}\mathrm{\Omega }_{mn,\widehat{\alpha }}+\gamma _{\widehat{\alpha }\widehat{\beta }}^p\mathrm{\Omega }_{mn,p}=0`$ \| \| $`(\frac{1}{2},\frac{3}{2},\frac{3}{2})`$ \| $`D_\alpha P^{\beta \widehat{\gamma }}\frac{1}{4}\left(\gamma ^{mn}\right)_\alpha ^\beta C_{mn}^{\widehat{\gamma }}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}P^{\beta \widehat{\gamma }}\frac{1}{4}\left(\gamma ^{pq}\right)_{\widehat{\alpha }}^{\widehat{\gamma }}C_{pq}^\beta =0`$ \| \| $`(\frac{1}{2},1,2)`$ \| $`D_\alpha \mathrm{\Omega }_{m,pq}_m\mathrm{\Omega }_{\alpha ,pq}\gamma _{m\alpha \gamma }C_{pq}^\gamma =0`$ \| $`\widehat{D}_{\widehat{\alpha }}\mathrm{\Omega }_{mn,p}+_p\mathrm{\Omega }_{mn,\widehat{\alpha }}+\gamma _{p\widehat{\alpha }\widehat{\beta }}C_{mn}^{\widehat{\beta }}=0`$ \| \| $`(\frac{1}{2},\frac{3}{2},2)`$ \| $`D_\alpha C_{pq}^\beta \frac{1}{4}\left(\gamma ^{mn}\right)_\alpha ^\beta S_{mn,pq}=0`$ \| $`\widehat{D}_{\widehat{\alpha }}C_{mn}^{\widehat{\beta }}+\frac{1}{4}\left(\gamma ^{pq}\right)_{\widehat{\alpha }}^{\widehat{\beta }}S_{mnpq}=0`$ \| (3.94) where the labels $`(a,b,c)`$ denote the scaling dimensions of the generators of the extended super-Poincaré algebra which the equations belong to. Moreover, one obtains the following eight equations, which do not provide further information, since they are implied by (3.94) and pure spinor conditions $`N^{mn}\lambda ^\gamma D_\gamma \mathrm{\Omega }_{mn,\widehat{\beta }}=0\widehat{\lambda }^{\widehat{\gamma }}\widehat{D}_{\widehat{\gamma }}\mathrm{\Omega }_{\alpha ,mn}\widehat{N}^{mn}=0`$ (3.95) $`N^{mn}\lambda ^\gamma D_\gamma \mathrm{\Omega }_{mn,p}=0\widehat{\lambda }^{\widehat{\gamma }}\widehat{D}_{\widehat{\gamma }}\mathrm{\Omega }_{m,pq}\widehat{N}^{pq}=0`$ (3.96) $`N^{mn}\lambda ^\gamma D_\gamma C_{mn}^{\widehat{\beta }}=0\widehat{\lambda }^{\widehat{\gamma }}\widehat{D}_{\widehat{\gamma }}C_{mn}^\alpha \widehat{N}^{mn}=0`$ (3.97) $`N^{mn}\lambda ^\gamma D_\gamma S_{mn,pq}\widehat{N}^{pq}=0N^{mn}\widehat{\lambda }^{\widehat{\gamma }}\widehat{D}_{\widehat{\gamma }}S_{mn,pq}\widehat{N}^{pq}=0`$ (3.98) Since we assumed that the superfields $`\mathrm{\Omega }_{mn,p},\mathrm{\Omega }_{m,pq},C_{mn}^{\widehat{\beta }},C_{pq}^\alpha `$ and $`S_{mn,pq}`$ correspond to the linearized curvatures of the connections, we can derive new equations needed for the iterative procedure outlined in the introduction. By contracting equations (3.94) with respect to the bosonic derivative and antisymmetrizing the bosonic indices, one obtains $`D_\alpha \mathrm{\Omega }_{mn,\widehat{\beta }}=_{[m}\gamma _{n]\alpha \gamma }E_{\widehat{\beta }}^\gamma \widehat{D}_{\widehat{\beta }}\mathrm{\Omega }_{\alpha ,pq}=_{[p}\gamma _{q]\widehat{\beta }\widehat{\gamma }}E_\alpha ^{\widehat{\gamma }}`$ (3.99) $`D_\alpha \mathrm{\Omega }_{mn,p}=_{[m}\gamma _{n]\alpha \gamma }E_p^\gamma \widehat{D}_{\widehat{\beta }}\mathrm{\Omega }_{m,pq}=_{[p}\gamma _{q]\widehat{\beta }\widehat{\gamma }}E_m^{\widehat{\gamma }}`$ (3.100) $`D_\alpha C_{mn}^{\widehat{\beta }}=_{[m}\gamma _{n]\alpha \gamma }P^{\gamma \widehat{\beta }}\widehat{D}_{\widehat{\beta }}C_{pq}^\alpha =_{[p}\gamma _{q]\widehat{\beta }\widehat{\gamma }}P^{\alpha \widehat{\gamma }}`$ (3.101) $`D_\alpha S_{mn,pq}=_{[m}\gamma _{n]\alpha \gamma }C_{pq}^\gamma \widehat{D}_{\widehat{\beta }}S_{mn,pq}=_{[p}\gamma _{q]\widehat{\beta }\widehat{\gamma }}C_{mn}^{\widehat{\gamma }}`$ (3.102) (we define $`a_{[m}b_{n]}=a_mb_na_nb_m`$). The identification of the superfields $`\mathrm{\Omega }_{mn,p}`$, $`\mathrm{\Omega }_{m,pq}`$, $`C_{mn}^{\widehat{\beta }}`$, $`C_{pq}^\alpha `$ and $`S_{mn,pq}`$ with the linearized curvatures is automatically derived in the formalism , and equations (3.99) are the usual Bianchi identities. In order to show that the above equations imply the supergravity equations of motion we proceed as follows. We first consider the third line of (3.99) and the $`(\frac{1}{2},\frac{3}{2},\frac{3}{2})`$ line of (3.94), that we recall for the reader convenience $`D_\alpha P^{\beta \widehat{\gamma }}{\displaystyle \frac{1}{4}}(\gamma ^{mn})_\alpha ^\beta C_{mn}^{\widehat{\gamma }}=0\widehat{D}_{\widehat{\alpha }}P^{\beta \widehat{\gamma }}{\displaystyle \frac{1}{4}}(\gamma ^{pq})_{\widehat{\alpha }}^{\widehat{\gamma }}C_{pq}^\beta =0`$ (3.104) $`D_\alpha C_{mn}^{\widehat{\beta }}=_{[m}\gamma _{n]\alpha \gamma }P^{\gamma \widehat{\beta }}\widehat{D}_{\widehat{\beta }}C_{pq}^\alpha =_{[p}\gamma _{q]\widehat{\beta }\widehat{\gamma }}P^{\alpha \widehat{\gamma }}`$ Acting with $`\gamma _{\alpha \sigma }^m_m`$ on $`P^{\sigma \widehat{\beta }}`$ and using the commutation relations of the $`D`$’s, one gets $`\gamma _{\alpha \sigma }^m_mP^{\sigma \widehat{\beta }}=(D_\alpha D_\sigma +D_\sigma D_\alpha )P^{\sigma \widehat{\beta }}={\displaystyle \frac{1}{4}}(\gamma ^{mn})_\alpha ^\sigma D_\sigma C_{mn}^{\widehat{\beta }},`$ (3.106) $`={\displaystyle \frac{1}{2}}(\gamma ^{mn})_\alpha ^\sigma \gamma _{m\sigma \gamma }_nP^{\gamma \widehat{\beta }}={\displaystyle \frac{9}{2}}\gamma _{\alpha \sigma }^m_mP^{\sigma \widehat{\beta }}`$ (3.107) Here we also used the first equation of (3.104) and $`D_\alpha P^{\alpha \widehat{\beta }}=0`$ (which follows from (3.104)). In the second line we used the first equation in the second line on (3.104) and the identity $`(\gamma ^{mn}\gamma _m)_{\alpha \beta }=9\gamma _{\alpha \beta }^n`$. By performing the same manipulations on the hatted quantities we derive the equations $$\gamma _{\alpha \sigma }^m_mP^{\sigma \widehat{\beta }}=0,\gamma _{\widehat{\alpha }\widehat{\sigma }}^m_mP^{\alpha \widehat{\sigma }}=0.$$ (3.108) Decomposing $`P^{\alpha \widehat{\beta }}`$ in terms of Dirac matrices, it is straightforward to show that (3.108) implies the equations of motion for the RR fields. Acting again with $`\gamma _n^{\alpha \gamma }D_\alpha `$ on (3.108) and using equations (3.104) one gets $$0=\gamma _n^{\alpha \gamma }\gamma _{\alpha \beta }^m_mD_\gamma P^{\beta \widehat{\beta }}=(\gamma _n\gamma ^m)_\beta ^\gamma (\gamma ^{pq})_\gamma ^\beta _mC_{pq}^{\widehat{\beta }}=^mC_{mn}^{\widehat{\beta }},$$ (3.109) and analogously for $`C_{pq}^\alpha `$. These equations are the Maxwell equations for the curvature of the gravitinos. They are not enough to describe the dynamic of gravitinos and we have to invoke new equations coming from the second line of (3.99) and the $`(\frac{1}{2},\frac{3}{2},1)`$ line of (3.99). Applying $`\gamma _{\alpha \sigma }^m_m`$ on $`E_p^\sigma `$ and with $`\gamma _{\widehat{\alpha }\widehat{\sigma }}^p_p`$ on $`E_m^{\widehat{\sigma }}`$, the same algebraic manipulations yield $$\gamma _{\alpha \sigma }^m_mE_p^\sigma =0,\gamma _{\widehat{\alpha }\widehat{\sigma }}^p_pE_m^{\widehat{\sigma }}=0.$$ (3.110) which are the Dirac equations for the gravitinos. These equations are gauge invariant under the gauge transformations discussed in the next section since the gauge parameters have to satisfy a field equation. In addition, as above, we find the equations $$^m\mathrm{\Omega }_{mn,p}=0,^p\mathrm{\Omega }_{m,pq}=0,$$ (3.111) which are, at the lowest component of the superfield $`\mathrm{\Omega }_{mn,p}`$ and $`\mathrm{\Omega }_{m,pq}`$, the equations of motion of the graviton, the dilaton and the NS-NS form $`^m(_{[m}g_{n]p}+_{[m}b_{n]p}+\eta _{p[n}_{m]}\varphi )=0,`$ (3.112) $`^p(_{[p}g_{\|m\|q]}+_{[p}b_{\|m\|q]}+\eta _{nm[q}_{p]}\varphi )=0.`$ (3.113) Pursuing this line of reasoning, one can derive similar equations for $`E_{\widehat{\alpha }}^\beta `$, $`E_\alpha ^{\widehat{\beta }}`$, $`\mathrm{\Omega }_{mn,\widehat{\gamma }}`$ and $`\mathrm{\Omega }_{pq}^\alpha `$, which guarantee that the fields are either pure gauge or auxiliary fields. Finally, by studying the last line of (3.99) and the line $`(\frac{1}{2},\frac{3}{2},2)`$ of (3.94), one derives new equations for $`C_{pq}^\beta ,C_{mn}^{\widehat{\beta }}`$ and $`S_{mn,pq}`$, which do not give further information since they are implied by the previous ones. #### 3.2.3 Gauge transformations and gauge fixing In order to solve the equations of motion (3.94) and (3.99) it is convenient to choose a suitable gauge. Indeed, for supersymmetric theories, the large amount of auxiliary fields can be reduced by choosing the Wess-Zumino gauge. We first discuss the general structure of the gauge transformations (3.58), we then provide a gauge fixing and we finally check that this gauge can be reached. In the present framework, the gauge parameters $`\mathrm{\Lambda }^{(0,1)}`$, $`\mathrm{\Lambda }^{(1,0)}`$, $`\tau _z^{(0,0)}`$ and $`\tau _{\overline{z}}^{(0,0)}`$ satisfy equations (3.60) and (3.61) and they are defined up to the gauge transformation (3.62). This additional gauge invariance is fixed by a further gauge fixing. The general structure of the gauge parameters $`\mathrm{\Lambda }^{(0,1)}`$, $`\mathrm{\Lambda }^{(1,0)}`$, $`\tau _z^{(0,0)}`$ and $`\tau _{\overline{z}}^{(0,0)}`$ is given by $`\mathrm{\Lambda }^{(1,0)}=\lambda ^\alpha \mathrm{\Theta }_\alpha \mathrm{\Lambda }^{(0,1)}=\widehat{\mathrm{\Theta }}_{\widehat{\alpha }}\widehat{\lambda }^{\widehat{\alpha }},`$ (3.114) and $`\tau _z^{(0,0)}=_z\theta ^\alpha \mathrm{\Xi }_\alpha +\mathrm{\Pi }_z^m\mathrm{\Sigma }_m+d_{z\alpha }\mathrm{\Phi }^\alpha +{\displaystyle \frac{1}{2}}N_z^{mn}\mathrm{\Psi }_{mn}`$ (3.115) $`\tau _{\overline{z}}^{(0,0)}=\widehat{\mathrm{\Xi }}_{\widehat{\alpha }}_{\overline{z}}\widehat{\theta }^{\widehat{\alpha }}+\widehat{\mathrm{\Sigma }}_p\widehat{\mathrm{\Pi }}_{\overline{z}}^p+\widehat{\mathrm{\Phi }}^{\widehat{\alpha }}\widehat{d}_{\overline{z}\widehat{\alpha }}+{\displaystyle \frac{1}{2}}\widehat{\mathrm{\Psi }}_{pq}\widehat{N}_{\overline{z}}^{pq}.`$ (3.116) where $`\mathrm{\Theta }_\alpha ,\mathrm{},\widehat{\mathrm{\Psi }}_{mn}`$ are superfields in the variables $`x^m,\theta ^\alpha `$, and $`\widehat{\theta }^{\widehat{\alpha }}`$. In terms of these superfields, eq. (3.60) gives $$(\gamma ^{mnpqr})^{\alpha \beta }D_\beta \mathrm{\Theta }_\alpha =0(\gamma ^{mnpqr})^{\widehat{\alpha }\widehat{\beta }}\widehat{D}_{\widehat{\beta }}\widehat{\mathrm{\Theta }}_{\widehat{\alpha }}=0,$$ (3.117) while eq. (3.61) gives $`\mathrm{\Theta }_\alpha +\mathrm{\Xi }_\alpha =0\widehat{\mathrm{\Theta }}_{\widehat{\alpha }}\widehat{\mathrm{\Xi }}_{\widehat{\alpha }}=0`$ (3.118) $`D_\alpha \mathrm{\Theta }_\beta D_\beta \mathrm{\Xi }_\alpha +\gamma _{\alpha \beta }^m\mathrm{\Sigma }_m=0\widehat{D}_{\widehat{\alpha }}\widehat{\mathrm{\Theta }}_{\widehat{\beta }}+\widehat{D}_{\widehat{\beta }}\widehat{\mathrm{\Xi }}_{\widehat{\alpha }}+\gamma _{\widehat{\alpha }\widehat{\beta }}^p\widehat{\mathrm{\Sigma }}_p=0`$ (3.119) $`D_\alpha \mathrm{\Sigma }_m+_m\mathrm{\Theta }_\alpha \gamma _{m\alpha \beta }\mathrm{\Phi }^\beta =0\widehat{D}_{\widehat{\alpha }}\widehat{\mathrm{\Sigma }}_p+_p\widehat{\mathrm{\Theta }}_{\widehat{\alpha }}+\gamma _{p\widehat{\alpha }\widehat{\beta }}\widehat{\mathrm{\Phi }}^{\widehat{\beta }}=0`$ (3.120) $`D_\alpha \mathrm{\Phi }^\beta {\displaystyle \frac{1}{4}}(\gamma ^{mn})_\alpha ^\beta \mathrm{\Psi }_{mn}=0\widehat{D}_{\widehat{\alpha }}\widehat{\mathrm{\Phi }}^{\widehat{\beta }}+{\displaystyle \frac{1}{4}}\left(\gamma ^{pq}\right)_{\widehat{\alpha }}^{\widehat{\beta }}\widehat{\mathrm{\Psi }}_{pq}`$ (3.121) $`N^{mn}\lambda ^\gamma D_\gamma \mathrm{\Psi }_{mn}=0\widehat{\lambda }^{\widehat{\gamma }}\widehat{D}_{\widehat{\gamma }}\widehat{\mathrm{\Psi }}_{pq}\widehat{N}^{pq}=0.`$ (3.122) These equations look like the superspace field equations for superMaxwell theory (3.29), however the superfields $`\mathrm{\Theta }_\alpha `$, $`\mathrm{\Sigma }_m`$, $`\mathrm{\Phi }^\alpha `$ and $`\mathrm{\Psi }_{mn}`$ and the corresponding hatted quantities depend on $`x^m`$, $`\theta ^\alpha `$ and $`\widehat{\theta }^{\widehat{\alpha }}`$. Therefore, the eqs. (3.118) are not sufficient to determine completely the components of those superfields. The free independent components are indeed the gauge parameters. We also note that the last pair of equations is trivial when the previous equations and the pure spinor conditions are imposed. Finally, because of the similarity with SYM case, it is quite natural to impose the condition that $`\mathrm{\Psi }_{mn}`$ and $`\widehat{\mathrm{\Psi }}_{pq}`$ are the linearized curvatures of $`\mathrm{\Sigma }_m`$ and $`\widehat{\mathrm{\Sigma }}_p`$. Again, this assumption is automatic in . The gauge transformations of the superfields in $`𝒱_{z\overline{z}}^{(0,0)}`$ are given by \| $`(\frac{1}{2},\frac{1}{2})`$ \| $`\delta A_{\alpha \widehat{\beta }}=D_\alpha \widehat{\mathrm{\Xi }}_{\widehat{\beta }}+\widehat{D}_{\widehat{\beta }}\mathrm{\Theta }_\alpha `$ \| \| \| --- \| --- \| --- \| \| $`(\frac{1}{2},1)`$ \| $`\delta A_{\alpha p}=_p\mathrm{\Theta }_\alpha +D_\alpha \widehat{\mathrm{\Sigma }}_p`$ \| $`\delta A_{m\widehat{\beta }}=_m\widehat{\mathrm{\Theta }}_{\widehat{\beta }}+\widehat{D}_{\widehat{\beta }}\mathrm{\Sigma }_m`$ \| \| $`(1,1)`$ \| $`\delta A_{mp}=_m\widehat{\mathrm{\Sigma }}_p_p\mathrm{\Sigma }_m`$ \| \| \| $`(\frac{3}{2},\frac{1}{2})`$ \| $`\delta E_{\widehat{\beta }}^\alpha =\widehat{D}_{\widehat{\beta }}\mathrm{\Phi }^\alpha `$ \| $`\delta E_\alpha ^{\widehat{\beta }}=D_\alpha \widehat{\mathrm{\Phi }}^{\widehat{\beta }}`$ \| \| $`(\frac{3}{2},1)`$ \| $`\delta E_p^\alpha =_p\mathrm{\Phi }^\alpha `$ \| $`\delta E_m^{\widehat{\beta }}=_m\widehat{\mathrm{\Phi }}^{\widehat{\beta }}`$ \| \| $`(\frac{1}{2},2)`$ \| $`\delta \mathrm{\Omega }_{\alpha ,pq}=D_\alpha \widehat{\mathrm{\Psi }}_{pq}`$ \| $`\delta \mathrm{\Omega }_{mn,\widehat{\beta }}=\widehat{D}_{\widehat{\beta }}\mathrm{\Psi }_{mn}`$ \| \| $`(1,2)`$ \| $`\delta \mathrm{\Omega }_{m,pq}=_m\widehat{\mathrm{\Psi }}_{pq}`$ \| $`\delta \mathrm{\Omega }_{mn,p}=_p\mathrm{\Psi }_{mn}`$ \| \| $`(\frac{3}{2},\frac{3}{2})`$ \| $`\delta P^{\alpha \widehat{\beta }}=0`$ \| \| \| $`(\frac{3}{2},2)`$ \| $`\delta C_{pq}^\alpha =0`$ \| $`\delta C_{mn}^{\widehat{\beta }}=0`$ \| \| $`(2,2)`$ \| $`\delta S_{mn,pq}=0`$ . \| \| (3.123) From these equations, we easily see that the superfields $`P^{\alpha \widehat{\beta }}`$, $`C_{pq}^\alpha `$, $`C_{mn}^{\widehat{\beta }}`$ and $`S_{mn,pq}`$ are indeed gauge invariant, as expected, being linearized field strengths. At zero order in $`\theta `$ and $`\widehat{\theta }`$ eq. (3.123) gives the gauge transformations of supergravity fields. For example, the first components of $`\widehat{\mathrm{\Sigma }}_p=\zeta _p+\xi _p+𝒪(\theta ,\widehat{\theta })`$ and $`\mathrm{\Sigma }_m=\zeta _m\xi _m+𝒪(\theta ,\widehat{\theta })`$ are to be identified with the parameters of diffeomorphisms $`\delta g_{mp}=_m\xi _p+_p\xi _m`$ and with the gauge transformations of the NS-NS form $`\delta b_{mp}=_m\zeta _p_p\zeta _m`$. So, the zero-order terms of the gauge parameter superfields $`\mathrm{\Theta }_\alpha `$, $`\mathrm{\Sigma }_m`$, $`\mathrm{\Phi }^\alpha `$ and of the corresponding hatted quantities are $`\mathrm{\Theta }_\alpha =𝒪(\theta ,\widehat{\theta });\widehat{\mathrm{\Theta }}_{\widehat{\beta }}=𝒪(\theta ,\widehat{\theta })`$ (3.124) $`\mathrm{\Sigma }_m=\zeta _m\xi _m+𝒪(\theta ,\widehat{\theta });\widehat{\mathrm{\Sigma }}_p=\zeta _p+\xi _p+𝒪(\theta ,\widehat{\theta })`$ (3.125) $`\mathrm{\Phi }^\alpha =\phi ^\alpha +𝒪(\theta ,\widehat{\theta });\widehat{\mathrm{\Phi }}^{\widehat{\beta }}=\widehat{\phi }^{\widehat{\beta }}+𝒪(\theta ,\widehat{\theta })`$ (3.126) Furthermore, the large amount of gauge parameters allows us to choose the gauge $`\theta ^\alpha A_{\alpha \widehat{\beta }}=0A_{\alpha \widehat{\beta }}\widehat{\theta }^{\widehat{\beta }}=0`$ (3.127) $`\theta ^\alpha A_{\alpha p}=0A_{m\widehat{\beta }}\widehat{\theta }^{\widehat{\beta }}=0`$ (3.128) $`\theta ^\alpha E_\alpha ^{\widehat{\beta }}=0E_{\widehat{\beta }}^\alpha \widehat{\theta }^{\widehat{\beta }}=0`$ (3.129) $`\theta ^\alpha \mathrm{\Omega }_{\alpha ,pq}=0\mathrm{\Omega }_{mn,\widehat{\beta }}\widehat{\theta }^{\widehat{\beta }}=0.`$ (3.130) Indeed, we have at our disposal the parameters $`\mathrm{\Theta }_\alpha `$, $`\mathrm{\Sigma }_m`$, $`\mathrm{\Phi }^\alpha `$ and $`\mathrm{\Psi }_{mn}`$ and the corresponding hatted quantities to impose the gauge (3.127). Before showing that the gauge can be reached we have to notice that the transformations (3.123) and the equations (3.117) are invariant under the residual gauge transformations (3.62) $`\delta \mathrm{\Theta }_\alpha =D_\alpha \mathrm{\Omega }\delta \widehat{\mathrm{\Xi }}_{\widehat{\beta }}=\widehat{D}_{\widehat{\beta }}\mathrm{\Omega }`$ (3.131) $`\delta \mathrm{\Sigma }_m=_m\mathrm{\Omega }\delta \widehat{\mathrm{\Sigma }}_p=_p\mathrm{\Omega }`$ (3.132) $`\delta \mathrm{\Phi }^\alpha =0\delta \widehat{\mathrm{\Phi }}^{\widehat{\beta }}=0`$ (3.133) $`\delta \mathrm{\Psi }_{mn}=0\delta \widehat{\mathrm{\Psi }}_{pq}=0,`$ (3.134) depending on the scalar superfield $`\mathrm{{\rm Y}}^{(0,0)}=\widehat{\mathrm{{\rm Y}}}^{(0,0)}\mathrm{\Omega }`$. This requires an additional gauge fixing $$\theta ^\alpha \mathrm{\Theta }_\alpha +\widehat{\theta }^{\widehat{\beta }}\widehat{\mathrm{\Xi }}_{\widehat{\beta }}=0.$$ (3.135) To show that the gauge choice (3.127) can be reached by the gauge transformations (3.123), we have to solve, for instance, the equations $$\theta ^\alpha (A_{\alpha \widehat{\beta }}+\delta A_{\alpha \widehat{\beta }})=0,(A_{\alpha \widehat{\beta }}+\delta A_{\alpha \widehat{\beta }})\widehat{\theta }^{\widehat{\beta }}=0,$$ (3.136) and analogously for all other gauge conditions (3.127). By using the properties of the superderivative, gauge fixing (3.135), consistency conditions (3.118), and by defining the operators $$𝐃\theta ^\alpha D_\alpha =\theta ^\alpha \frac{}{\theta ^\alpha },\widehat{𝐃}\widehat{\theta }^{\widehat{\beta }}\widehat{D}_{\widehat{\beta }}=\widehat{\theta }^{\widehat{\beta }}\frac{}{\widehat{\theta }^{\widehat{\beta }}},$$ (3.137) we get the following recursive equations \| $`\left(1+𝐃+\widehat{𝐃}\right)\mathrm{\Theta }_\alpha =A_{\alpha \widehat{\beta }}\widehat{\theta }^{\widehat{\beta }}\left(\gamma ^m\theta \right)_\alpha \mathrm{\Sigma }_m`$ \| $`\left(1+𝐃+\widehat{𝐃}\right)\widehat{\mathrm{\Theta }}_{\widehat{\beta }}=\theta ^\alpha A_{\alpha \widehat{\beta }}\left(\gamma ^p\widehat{\theta }\right)_{\widehat{\beta }}\widehat{\mathrm{\Sigma }}_p`$ \| \| --- \| --- \| \| $`\left(𝐃+\widehat{𝐃}\right)\mathrm{\Sigma }_m=A_{m\widehat{\beta }}\widehat{\theta }^{\widehat{\beta }}+\left(\gamma _m\theta \right)_\beta \mathrm{\Phi }^\beta `$ \| $`\left(𝐃+\widehat{𝐃}\right)\widehat{\mathrm{\Sigma }}_p=\theta ^\alpha A_{\alpha p}\left(\gamma _p\widehat{\theta }\right)_{\widehat{\gamma }}\widehat{\mathrm{\Phi }}^{\widehat{\gamma }}`$ \| \| $`\left(𝐃+\widehat{𝐃}\right)\mathrm{\Phi }^\alpha =E_{\widehat{\beta }}^\alpha \widehat{\theta }^{\widehat{\beta }}\frac{1}{4}\left(\gamma ^{mn}\theta \right)^\alpha \mathrm{\Psi }_{mn}`$ \| $`\left(𝐃+\widehat{𝐃}\right)\widehat{\mathrm{\Phi }}^{\widehat{\beta }}=\theta ^\alpha E_\alpha ^{\widehat{\beta }}+\frac{1}{4}\left(\gamma ^{pq}\widehat{\theta }\right)^\beta \widehat{\mathrm{\Psi }}_{pq}`$ \| \| $`\left(𝐃+\widehat{𝐃}\right)\mathrm{\Psi }_{mn}=\mathrm{\Omega }_{mn,\widehat{\beta }}\widehat{\theta }^{\widehat{\beta }}\left(\gamma _{[m}\theta \right)_\gamma _{n]}\mathrm{\Phi }^\gamma `$ \| $`\left(𝐃+\widehat{𝐃}\right)\widehat{\mathrm{\Psi }}_{pq}=\theta ^\alpha \mathrm{\Omega }_{\alpha ,pq}+\left(\gamma _{[p}\widehat{\theta }\right)_{\widehat{\gamma }}_{q]}\widehat{\mathrm{\Phi }}^{\widehat{\gamma }}`$ \| (3.138) The operator $`(𝐃+\widehat{𝐃})`$ acts on homogeneous polynomials in $`\theta ^\alpha `$ and $`\widehat{\theta }^{\widehat{\alpha }}`$ by multiplication by the degree of homogeneity and it does not change its degree. Therefore, the relations (3.138) are recursive in powers of $`\theta `$ and $`\widehat{\theta }`$. They can be solved algebraically given $`A_{\alpha \widehat{\beta }},\mathrm{},\mathrm{\Omega }_{\alpha ,pq}`$ order by order in $`\theta `$ and $`\widehat{\theta }`$ and this proves that the gauge can indeed be imposed. Of course, to reconstruct the gauge-parameter superfields by means of the recursive equations (3.138), we also need lowest order data for them. These are the zero order supergravity gauge parameters (3.124). To obtain the last couple of equations we used the additional condition that $`\mathrm{\Psi }_{mn}`$ and $`\widehat{\mathrm{\Psi }}_{pq}`$ are the linearized curvatures of $`\mathrm{\Sigma }_m`$ and $`\widehat{\mathrm{\Sigma }}_p`$. #### 3.2.4 Recursive equations and explicit solution (up to order $`\theta ^2\widehat{\theta }^2`$) The next step is the derivation of the recursion equations for supergravity superfields. Acting with $`D_\alpha `$ and $`\widehat{D}_{\widehat{\beta }}`$ on the gauge fixing conditions (3.127), and using the definition (3.137), it is straightforward to derive the recursion relations from eq. (3.94) \| $`(1+𝐃)A_{\alpha \widehat{\beta }}=(\gamma ^m\theta )_\alpha A_{m\widehat{\beta }}`$ \| $`(1+\widehat{𝐃})A_{\alpha \widehat{\beta }}=(\gamma ^p\widehat{\theta })_{\widehat{\beta }}A_{\alpha p}`$ \| \| --- \| --- \| \| $`𝐃A_{m\widehat{\beta }}=(\gamma _m\theta )_\gamma E_{\widehat{\beta }}^\gamma `$ \| $`\widehat{𝐃}A_{\alpha p}=(\gamma _p\widehat{\theta })_{\widehat{\gamma }}E_\alpha ^{\widehat{\gamma }}`$ \| \| $`𝐃E_{\widehat{\beta }}^\alpha =\frac{1}{4}(\gamma ^{mn}\theta )^\alpha \mathrm{\Omega }_{mn,\widehat{\beta }}`$ \| $`\widehat{𝐃}E_\alpha ^{\widehat{\beta }}=\frac{1}{4}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}\mathrm{\Omega }_{\alpha ,pq}`$ \| \| $`𝐃\mathrm{\Omega }_{mn,\widehat{\beta }}=(\gamma _{[m}\theta )_\gamma _{n]}E_{\widehat{\beta }}^\gamma `$ \| $`\widehat{𝐃}\mathrm{\Omega }_{\alpha ,pq}=(\gamma _{[p}\widehat{\theta })_{\widehat{\gamma }}_{q]}E_\alpha ^{\widehat{\gamma }}`$ \| (3.139) \| $`(1+𝐃)A_{\alpha p}=(\gamma ^m\theta )_\alpha A_{mp}`$ \| $`(1+\widehat{𝐃})A_{m\widehat{\beta }}=(\gamma ^p\widehat{\theta })_{\widehat{\beta }}A_{mp}`$ \| \| --- \| --- \| \| $`𝐃A_{mp}=(\gamma _m\theta )_\beta E_p^\beta `$ \| $`\widehat{𝐃}A_{mp}=(\gamma _p\widehat{\theta })_{\widehat{\beta }}E_m^{\widehat{\beta }}`$ \| \| $`𝐃E_p^\alpha =\frac{1}{4}(\gamma ^{mn}\theta )^\alpha \mathrm{\Omega }_{mn,p}`$ \| $`\widehat{𝐃}E_m^{\widehat{\beta }}=\frac{1}{4}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}\mathrm{\Omega }_{m,pq}`$ \| \| $`𝐃\mathrm{\Omega }_{mn,p}=(\gamma _{[m}\theta )_\gamma _{n]}E_p^\gamma `$ \| $`\widehat{𝐃}\mathrm{\Omega }_{m,pq}=(\gamma _{[p}\widehat{\theta })_{\widehat{\gamma }}_{q]}E_m^{\widehat{\gamma }}`$ \| (3.140) \| $`(1+𝐃)E_\alpha ^{\widehat{\beta }}=(\gamma ^m\theta )_\alpha E_m^{\widehat{\beta }}`$ \| $`(1+\widehat{𝐃})E_{\widehat{\beta }}^\alpha =(\gamma ^p\widehat{\theta })_{\widehat{\beta }}E_p^\alpha `$ \| \| --- \| --- \| \| $`𝐃E_m^{\widehat{\beta }}=(\gamma _m\theta )_\gamma P^{\gamma \widehat{\beta }}`$ \| $`\widehat{𝐃}E_p^\alpha =(\gamma _p\widehat{\theta })_{\widehat{\gamma }}P^{\alpha \widehat{\gamma }}`$ \| \| $`𝐃P^{\alpha \widehat{\beta }}=\frac{1}{4}(\gamma ^{mn}\theta )^\alpha C_{mn}^{\widehat{\beta }}`$ \| $`\widehat{𝐃}P^{\alpha \widehat{\beta }}=\frac{1}{4}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}C_{pq}^\alpha `$ \| \| $`𝐃C_{mn}^{\widehat{\beta }}=(\gamma _{[m}\theta )_\gamma _{n]}P^{\gamma \widehat{\beta }}`$ \| $`\widehat{𝐃}C_{pq}^\alpha =(\gamma _{[p}\widehat{\theta })_{\widehat{\gamma }}_{q]}P^{\alpha \widehat{\gamma }}`$ \| (3.141) \| $`(1+𝐃)\mathrm{\Omega }_{\alpha ,pq}=(\gamma ^m\theta )_\alpha \mathrm{\Omega }_{m,pq}`$ \| $`(1+\widehat{𝐃})\mathrm{\Omega }_{mn,\widehat{\beta }}=(\gamma ^p\widehat{\theta })_{\widehat{\beta }}\mathrm{\Omega }_{mn,p}`$ \| \| --- \| --- \| \| $`𝐃\mathrm{\Omega }_{m,pq}=(\gamma _m\theta )_\beta C_{pq}^\beta `$ \| $`\widehat{𝐃}\mathrm{\Omega }_{mn,p}=(\gamma _p\widehat{\theta })_{\widehat{\beta }}C_{mn}^{\widehat{\beta }}`$ \| \| $`𝐃C_{pq}^\alpha =\frac{1}{4}(\gamma ^{mn}\theta )^\alpha S_{mn,pq}`$ \| $`\widehat{𝐃}C_{mn}^{\widehat{\beta }}=\frac{1}{4}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}S_{mn,pq}`$ \| \| $`𝐃S_{mn,pq}=(\gamma _{[m}\theta )_\gamma _{n]}C_{pq}^\gamma ,`$ \| $`\widehat{𝐃}S_{mn,pq}=(\gamma _{[p}\widehat{\theta })_{\widehat{\gamma }}_{q]}C_{mn}^{\widehat{\gamma }}`$ \| (3.142) A given superfield appears in two groups of equations in order that both its $`\theta `$ and $`\widehat{\theta }`$ components are fixed. Inside each group there is an iterative structure (see and ) which allows us to solve those equations recursively given the initial conditions and there is a hierarchical structure among the different groups of equations which allows us to solve them subsequently. To provide the initial data, we identify the lowest-components of the matrix superfield $`𝐀`$ in (3.85) with supergravity fields $$𝐀=\left[\begin{array}{cccc}0& 0& 0& 0\\ 0& g_{mp}+b_{mp}+\eta _{mp}\varphi & \psi _m^{\widehat{\beta }}& \omega _{m,pq}\\ 0& \psi _p^\alpha & f^{\alpha \widehat{\beta }}& c_{pq}^\alpha \\ 0& \omega _{mn,p}& c_{mn}^{\widehat{\beta }}& s_{mn,pq}\end{array}\right]+𝒪(\theta ,\widehat{\theta }),$$ (3.143) where the linearized gravitational connection and curvatures are given by $`\omega _{m,pq}=(_pg_{mq}_qg_{mp})+(_pb_{mq}_qb_{mp})+(\eta _{mq}_p\eta _{mp}_q)\varphi ,`$ (3.144) $`c_{mn}^{\widehat{\beta }}=(_m\psi _n^{\widehat{\beta }}_n\psi _m^{\widehat{\beta }}),`$ (3.145) $`s_{mn,pq}=(_m\omega _{n,pq}_n\omega _{m,pq}),`$ (3.146) and, analogously, for $`\omega _{mn,p}`$ and $`c_{pq}^\alpha `$. In the following we give the component-expansion for the physical superfields $`A_{mp}`$, $`E_m^{\widehat{\beta }}`$, $`E_p^\alpha `$ and $`P^{\alpha \widehat{\beta }}`$, up to second order in both $`\theta `$ and $`\widehat{\theta }`$. The corresponding curvatures can be easily computed from the defining equations (3.144). $`A_{mp}`$ $`=(g+b+\eta \varphi )_{mp}+(\gamma _m\theta )_\beta \psi _p^\beta (\gamma _p\widehat{\theta })_{\widehat{\beta }}\psi _m^{\widehat{\beta }}+(\gamma _m\theta )_\beta (\gamma _p\widehat{\theta })_{\widehat{\gamma }}f^{\beta \widehat{\gamma }}`$ (3.150) $`{\displaystyle \frac{1}{8}}(\gamma _m\theta )_\beta (\gamma ^{nr}\theta )^\beta \omega _{nr,p}{\displaystyle \frac{1}{8}}(\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\beta }}\omega _{m,qr}`$ $`+{\displaystyle \frac{1}{8}}(\gamma _m\theta )_\beta (\gamma ^{nr}\theta )^\beta (\gamma _p\widehat{\theta })_{\widehat{\gamma }}c_{nr}^{\widehat{\gamma }}{\displaystyle \frac{1}{8}}(\gamma _m\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\beta }}c_{qr}^\gamma `$ $`+{\displaystyle \frac{1}{64}}(\gamma _m\theta )_\beta (\gamma ^{nr}\theta )^\beta (\gamma _p\widehat{\theta })_{\widehat{\gamma }}(\gamma ^{qs}\widehat{\theta })^{\widehat{\gamma }}s_{nr,qs}+\mathrm{}`$ $`E_m^{\widehat{\beta }}`$ $`=\psi _m^{\widehat{\beta }}+(\gamma _m\theta )_\gamma f^{\gamma \widehat{\beta }}+{\displaystyle \frac{1}{4}}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}\omega _{m,pq}{\displaystyle \frac{1}{4}}(\gamma _m\theta )_\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}c_{pq}^\gamma `$ (3.154) $`{\displaystyle \frac{1}{8}}(\gamma _m\theta )_\gamma (\gamma ^{nr}\theta )^\gamma c_{nr}^{\widehat{\beta }}+{\displaystyle \frac{1}{4}}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_q\psi _m^{\widehat{\gamma }}`$ $`{\displaystyle \frac{1}{32}}(\gamma _m\theta )_\gamma (\gamma ^{nr}\theta )^\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}s_{nr,pq}+{\displaystyle \frac{1}{4}}(\gamma _m\theta )_\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_qf^{\gamma \widehat{\gamma }}`$ $`{\displaystyle \frac{1}{32}}(\gamma _m\theta )_\gamma (\gamma ^{nr}\theta )^\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_qc_{nr}^{\widehat{\gamma }}+\mathrm{}`$ $`E_p^\alpha `$ $`=\psi _p^\alpha {\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha \omega _{mn,p}+(\gamma _p\widehat{\theta })_{\widehat{\gamma }}f^{\alpha \widehat{\gamma }}+{\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha (\gamma _p\widehat{\theta })_{\widehat{\beta }}c_{mn}^{\widehat{\beta }}`$ (3.158) $`+{\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma _n\psi _p^\gamma {\displaystyle \frac{1}{8}}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\gamma }}c_{qr}^\alpha `$ $`+{\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}_nf^{\gamma \widehat{\beta }}+{\displaystyle \frac{1}{32}}(\gamma ^{mn}\theta )^\alpha (\gamma _p\widehat{\theta })_{\widehat{\gamma }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\gamma }}s_{mn,qr}`$ $`+{\displaystyle \frac{1}{32}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\beta }}_nc_{qr}^\gamma +\mathrm{}`$ $`P^{\alpha \widehat{\beta }}`$ $`=f^{\alpha \widehat{\beta }}{\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha c_{mn}^{\widehat{\beta }}{\displaystyle \frac{1}{4}}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}c_{pq}^\alpha {\displaystyle \frac{1}{16}}(\gamma ^{mn}\theta )^\alpha (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}s_{mn,pq}`$ (3.162) $`+{\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma _nf^{\gamma \widehat{\beta }}+{\displaystyle \frac{1}{4}}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_qf^{\alpha \widehat{\gamma }}`$ $`+{\displaystyle \frac{1}{16}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}_nc_{pq}^\gamma {\displaystyle \frac{1}{16}}(\gamma ^{mn}\theta )^\alpha (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_qc_{mn}^{\widehat{\gamma }}`$ $`+{\displaystyle \frac{1}{16}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_n_qf^{\gamma \widehat{\gamma }}+\mathrm{}`$ In the following we list the solution up to second order in both $`\theta ^\alpha `$ and $`\widehat{\theta }^{\widehat{\alpha }}`$ for the auxiliary superfields $`A_{\alpha \widehat{\beta }},A_{\alpha p},A_{m\widehat{\beta }},E_\alpha ^{\widehat{\beta }}`$ and $`E_{\widehat{\beta }}^\alpha `$. $`A_{\alpha \widehat{\beta }}`$ $`={\displaystyle \frac{1}{4}}(\gamma ^m\theta )_\alpha (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(g+b+\eta \varphi )_{mp}`$ (3.165) $`+{\displaystyle \frac{1}{6}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}\psi _p^\gamma +{\displaystyle \frac{1}{6}}(\gamma ^m\theta )_\alpha (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}\psi _m^{\widehat{\gamma }}`$ $`+{\displaystyle \frac{1}{9}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\beta (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}f^{\beta \widehat{\gamma }}+\mathrm{}`$ $`A_{\alpha p}`$ $`={\displaystyle \frac{1}{2}}\theta ^\beta \gamma _{\beta \alpha }^m(g+b+\eta \varphi )_{mp}{\displaystyle \frac{1}{2}}(\gamma ^m\theta )_\alpha (\gamma _p\widehat{\theta })_{\widehat{\gamma }}\psi _m^{\widehat{\gamma }}+{\displaystyle \frac{1}{3}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\beta \psi _p^\beta `$ (3.169) $`+{\displaystyle \frac{1}{3}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\beta (\gamma _p\widehat{\theta })_{\widehat{\gamma }}f^{\beta \widehat{\gamma }}{\displaystyle \frac{1}{16}}(\gamma ^m\theta )_\alpha (\gamma _p\widehat{\theta })_{\widehat{\gamma }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\gamma }}\omega _{m,qr}`$ $`{\displaystyle \frac{1}{24}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\beta }}c_{qr}^\gamma +\mathrm{}`$ $`A_{m\widehat{\beta }}`$ $`={\displaystyle \frac{1}{2}}\widehat{\theta }^{\widehat{\gamma }}\gamma _{\widehat{\gamma }\widehat{\beta }}^p(g+b+\eta \varphi )_{mp}+{\displaystyle \frac{1}{2}}(\gamma _m\theta )_\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}\psi _p^\gamma +{\displaystyle \frac{1}{3}}(\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}\psi _m^{\widehat{\gamma }}`$ (3.172) $`+{\displaystyle \frac{1}{16}}(\gamma _m\theta )_\gamma (\gamma ^{nr}\theta )^\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}\omega _{nr,p}+{\displaystyle \frac{1}{3}}(\gamma _m\theta )_\beta (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}f^{\beta \widehat{\gamma }}`$ $`{\displaystyle \frac{1}{24}}(\gamma _m\theta )_\gamma (\gamma ^{nr}\theta )^\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}c_{nr}^{\widehat{\gamma }}+\mathrm{}`$ $`E_\alpha ^{\widehat{\beta }}`$ $`={\displaystyle \frac{1}{2}}\theta ^\gamma \gamma _{\gamma \alpha }^m\psi _m^{\widehat{\beta }}+{\displaystyle \frac{1}{8}}(\gamma ^m\theta )_\alpha (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}\omega _{m,pq}+{\displaystyle \frac{1}{3}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma f^{\gamma \widehat{\beta }}`$ (3.175) $`{\displaystyle \frac{1}{12}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}c_{pq}^\gamma +{\displaystyle \frac{1}{8}}(\gamma ^m\theta )_\alpha (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_q\psi _m^{\widehat{\gamma }}`$ $`+{\displaystyle \frac{1}{12}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_qf^{\gamma \widehat{\gamma }}+\mathrm{}`$ $`E_{\widehat{\beta }}^\alpha `$ $`={\displaystyle \frac{1}{2}}\widehat{\theta }^{\widehat{\gamma }}\gamma _{\widehat{\gamma }\widehat{\beta }}^p\psi _p^\alpha +{\displaystyle \frac{1}{8}}(\gamma ^{mn}\theta )^\alpha (\gamma ^p\widehat{\theta })_{\widehat{\beta }}\omega _{mn,p}+{\displaystyle \frac{1}{3}}(\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}f^{\alpha \widehat{\gamma }}`$ (3.178) $`{\displaystyle \frac{1}{8}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}_n\psi _p^\gamma {\displaystyle \frac{1}{12}}(\gamma ^{mn}\theta )^\alpha (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}c_{mn}^{\widehat{\gamma }}`$ $`+{\displaystyle \frac{1}{12}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}_nf^{\gamma \widehat{\gamma }}+\mathrm{}`$ It is easy to verify that this expansion satisfies the gauge conditions (3.127) and that all auxiliary fields have been eliminated and reexpressed in terms of derivatives of physical supergravity fields. The next step is to insert the expansion (3.154) into the definition of the vertex operator (3.89) and recombine the worldsheet one-forms $`𝐗_z`$ and $`𝐗_{\overline{z}}`$ in order to get a more manageable expression. However, it makes sense to provide such expression for an interesting example in section 3.3. We have to notice that the vertex operator $`𝒱^{(1,1)}`$ contains only the superfield $`A_{\alpha \widehat{\beta }}`$ which encodes all the needed information regarding the supergravity fields, which however appear at higher orders in $`\theta `$’s and $`\widehat{\theta }`$’s. This is sufficient for amplitudes computations, even though the measure factor on zero modes in the correlation functions has to soak up plenty of $`\theta `$’s and $`\widehat{\theta }`$’s (). #### 3.2.5 Gauge fixing for massive states In the previous sections, we explored the gauge fixing for the massless sector of open and closed string theory. However, the spectrum of string theory contains infinitely many massive states defined, in the closed string case, by the equations $$[Q_L,𝒱_n^{(1,1)}]=0,[Q_R,𝒱_n^{(1,1)}]=0,[L_{0,L}+L_{0,R}n,𝒱_n^{(1,1)}]=0,$$ (3.179) where $`L_{0,L}=𝑑zzT_{zz}`$ and $`L_{0,R}=𝑑\overline{z}\overline{z}\overline{T}_{\overline{z}\overline{z}}`$. The index $`n`$ denotes the mass of the state. Even if these equations can be solved by expanding the vertex operators $`𝒱_n^{(1,1)}`$ in terms of the building-blocks $`\theta ^\alpha `$, $`\overline{}\widehat{\theta }^{\widehat{\alpha }}`$, $`\mathrm{\Pi }^m`$, $`\overline{\mathrm{\Pi }}^m`$,…, it is convenient to fix a gauge as in the massless case and then solve the equations by an iterative construction as shown in the previous section. However, since we cannot explore the complete set of vertices and provide a gauge fixing for each of them, we propose a definition of gauge fixing based on new anticommuting and nilpotent charges to be imposed on the physical states. This resembles the Siegel gauge (where the corresponding charges are $`b_{L,0}=𝑑zzb_{zz}`$ and $`b_{L,0}=𝑑\overline{z}\overline{z}\widehat{b}_{\overline{z}\overline{z}}`$ where $`b_{zz}`$ and $`\widehat{b}_{\overline{z}\overline{z}}`$ are the left- and right-moving antighosts) used in string field theory to eliminate all auxiliary fields and to define the propagator for the string field. We introduce the following charges “dual” to the BRST operators $$𝒦_L=𝑑z\theta ^\alpha w_\alpha ,𝒦_R=𝑑\overline{z}\widehat{\theta }^{\widehat{\beta }}\widehat{w}_{\widehat{\beta }}.$$ (3.180) They are nilpotent and anti-commute. They are not supersymmetry invariant as can be directly seen by the presence of $`\theta ^\alpha `$ and $`\widehat{\theta }^{\widehat{\beta }}`$. This in fact implies that we are choosing a non symmetric gauge which can be viewed as a generalization of the Wess-Zumino gauge condition in 10 dimensions. It eliminates the lowest non physical component of the superfields and it fixes the auxiliary fields – appearing at higher order in the superspace expansion – in terms of the physical fields and their derivatives. In addition, $`𝒦_{L/R}`$ are not invariant under the gauge transformations (3.35) , but their gauge variations are BRST invariant because of the pure spinor conditions $$\{Q_L,\mathrm{\Delta }_L𝒦_L\}=0,\{Q_R,\mathrm{\Delta }_R𝒦_R\}=0,$$ (3.181) Moreover, $`𝒦_{L/R}`$ have the following commutation relations with the BRST operators $`\{Q_L,𝒦_L\}=𝒟+J_L,\{Q_R,𝒦_L\}=0,`$ (3.182) $`\{Q_R,𝒦_R\}=\widehat{𝒟}+J_R,\{Q_L,𝒦_R\}=0,`$ (3.183) where $`𝒟={\displaystyle }dz:\theta ^\alpha d_\alpha :,J_L={\displaystyle }dz:\lambda ^\alpha w_\alpha :`$ (3.184) $`\widehat{𝒟}={\displaystyle }d\overline{z}:\widehat{\theta }^{\widehat{\alpha }}\widehat{d}_{\widehat{\alpha }}:,J_R={\displaystyle }d\overline{z}:\lambda ^{\widehat{\alpha }}\widehat{w}_{\widehat{\alpha }}:`$ (3.185) Acting on superfields $`F(x,\theta ,\widehat{\theta })`$, we have that $`\{𝒟,F\}=𝐃F`$ and $`\{\widehat{𝒟},F\}=\widehat{𝐃}F`$. The ordering of fields in the operators $`𝒟`$, $`\widehat{𝒟}`$, $`J_L`$ and $`J_R`$ is needed to define the corresponding currents. The operators are gauge invariant under (3.35) because of (3.181). The main difference with respect to Siegel gauge fixing in string field theory is that in that case $`b_{zz}`$ and $`\widehat{b}_{\widehat{z}\widehat{z}}`$ are holomorphic and antiholomorphic anticommuting currents of spin 2. In the case of the open superstring, denoting by $`Q`$ and by $`𝒦`$ the BRST and gauge fixing operators, the gauge condition on the massless vertex operator $`𝒱^{(1)}=\lambda ^\alpha A_\alpha `$ is given by $$\{𝒦,𝒱^{(1)}\}=𝑑w\left(\theta ^\alpha w_\alpha \right)(w)\left(\lambda ^\alpha A_\alpha (x,\theta )\right)(z)=\theta ^\alpha A_\alpha =0.$$ (3.186) We notice that the field $`\theta ^\alpha `$ in $`𝒦`$ is harmless for massless vertices, but it will give a nontrivial contribution in the massive case. In the latter case one has to add a compensating non-gauge invariant contribution on the r.h.s. of (3.186) in order to compensate the fact that $`𝒦`$ is not gauge invariant under (3.35). Applying $`Q`$ on the left hand side of (3.186) applying $`𝒦`$ on the equation $`\{Q,𝒱^{(1)}\}=\lambda \gamma ^m\lambda A_m(x,\theta )=0`$ and using the commutation relations (3.182), we obtain $$(𝐃+1)𝒱^{(1)}=\lambda \gamma ^m\theta A_m.$$ (3.187) Eliminating the ghost $`\lambda ^\alpha `$, we end up with equation (3.77) for the superfields $`A_\alpha `$ and $`A_m`$. This procedure can be clearly generalized to massive states. First, we discuss the closed string case, then we show an example for the first massive state for open superstrings and, finally, we show that the zero momentum cohomology satisfies the same equations generalized to zero modes. For closed strings, we reproduce the gauge fixing (3.127) by the following conditions $$\{𝒦_L,𝒱^{(1,1)}\}=0,\{𝒦_R,𝒱^{(1,1)}\}=0$$ (3.188) and, for the gauge parameters $`\mathrm{\Lambda }^{(1,0)}`$ and $`\mathrm{\Lambda }^{(0,1)}`$ in eq. (3.114), by the gauge condition $$\{𝒦_L,\mathrm{\Lambda }^{(1,0)}\}+\{𝒦_R,\mathrm{\Lambda }^{(0,1)}\}=0.$$ (3.189) which coincides with (3.135). Applying the BRST charge on the left hand sides of (3.188), acting with $`𝒦_L`$ and $`𝒦_R`$ on equations (3.51), and finally using the commutation relations (3.182), we derive the conditions for the iterative equations given in the previous section. Let us show that the gauge fixing (3.186) also fixes the gauge transformations in a suitable way for the first massive state of the open superstring $`𝒱_1^{(1)}`$, leading to a recursive procedure to compute the vertex operator in term of the initial data, a multiplet of on-shell fields containing a massive spin 2 field . A general decomposition of $`𝒱_1^{(1)}`$ in terms of fundamental building-blocks is given in (3.23) and its gauge transformation is generated by $$\delta 𝒱_1^{(1)}=[Q,\mathrm{\Omega }_1^{(0)}],$$ (3.190) with $$\mathrm{\Omega }_1^{(0)}=\theta ^\beta \mathrm{\Omega }_\beta +:\mathrm{\Pi }^m\mathrm{\Omega }_m:+:d_\beta \mathrm{\Omega }^\beta :+:N^{mn}:\mathrm{\Omega }_{mn}+:w_\beta \lambda ^\beta :\mathrm{\Omega }.$$ (3.191) The decompositions are based on the requirement that the vertex operator should be invariant under the gauge transformation $`\mathrm{\Delta }`$ given in (3.35). A further gauge transformation of $`\mathrm{\Omega }_1^{(0)}`$ would be a variation of a negative ghost number field. The only one is the antighost $`w_\alpha `$, but there is no gauge invariant operator only with $`w_\alpha `$ without $`\lambda ^\alpha `$. Notice that we have to add a (BRST-invariant) compensating term of the form $`w\gamma ^{mnpq}\lambda `$ in order to reabsorb the non-invariance of $`𝒦`$. Imposing (3.186), we get $`A_\alpha +\theta ^\beta B_{\beta \alpha }=0,\theta ^\alpha H_{\alpha m}=0,\theta ^\beta C_\beta ^\alpha =0,`$ (3.192) $`\theta ^\beta F_{\beta mn}+{\displaystyle \frac{1}{1440}}\left[(\gamma _{mn})_\gamma ^\alpha C_\alpha ^\gamma (\gamma _{mn})_\gamma ^\alpha \theta ^\gamma E_\alpha \right]=0,`$ (3.193) $`{\displaystyle \frac{1}{2}}(\gamma ^{mn}\theta )^\beta F_{\beta mn}C_\alpha ^\alpha +2\theta ^\alpha E_\alpha =0.`$ (3.194) This gauge fixing can be reached by adjusting the parameters $`\mathrm{\Omega }_\alpha ,\mathrm{\Omega }_m,\mathrm{\Omega }^\alpha ,\mathrm{\Omega }_{mn}`$ and $`\mathrm{\Omega }`$. Using equations (3.192) and applying the operator $`𝐃`$, we obtain the iterative relations to compute the vertex. The gauge fixing (3.192) fixes only the supergauge part of the gauge transformation. This gauge does not fix the physical gauge transformation of the massive spin 2 system . Finally, we show that the measure for zero modes satisfies the gauge fixing proposed above. In fact, by restricting the attention to zero momentum cohomology, we supersede $`𝒦`$ with the differential $$𝒦_0=\theta _0^\alpha \frac{}{\lambda _0^\alpha }$$ (3.195) which acting on $`𝒱^{(3)}=(\lambda _0\gamma ^m\theta _0)(\lambda _0\gamma ^n\theta _0)(\lambda _0\gamma ^p\theta _0)(\theta _0\gamma _{mnp}\theta _0)`$, yields $$𝒦_0𝒱^{(3)}=0.$$ (3.196) Similarly, for the closed superstring, the ghost number $`(3,3)`$ cohomology representative satisfies the corresponding gauge fixing. Even if the gauge fixing is not manifestly supersymmetric, the supersymmetry of the target space theory is still a symmetry. As usual, in the Wess-Zumino gauge, a supersymmetry transformation must be accompanied by a gauge transformation to bring the vertex to the original gauge. This means that $$\delta _ϵ[𝒦,𝒱]+[𝒦,\delta 𝒱]=0$$ (3.197) where $`\delta 𝒱=[Q,\mathrm{\Omega }_ϵ]`$, $`\delta _ϵ𝒱=[ϵ^\alpha Q_\alpha ,𝒱]`$ and $`Q_\alpha =𝑑zq_\alpha `$ (the supersymmetry generator $`q_\alpha `$ is given in (3.39)). As an example, we show that $`\mathrm{\Omega }_ϵ`$ can be indeed found for the massless sector of the open superstring and the extension is similar for the other cases. Equation (3.197) reduces to $$ϵ^\alpha A_\alpha +\theta ^\alpha ϵ^\beta Q_\beta A_\alpha +\theta ^\alpha D_\alpha \mathrm{\Omega }_ϵ=0,$$ (3.198) which yields $$𝐃\mathrm{\Omega }_ϵ=0.$$ (3.199) Again, this equation can be solved iteratively in powers of $`\theta `$’s and it follows that $`\mathrm{\Omega }=\mathrm{\Omega }_0(x)`$. (3.199) can be checked explicitly on the solutions (3.154). ### 3.3 An example: Linearly $`x`$-dependent Ramond-Ramond field strength and possible Lie-algebraic superspace deformations In this section I will discuss in detail one application of the iterative procedure I presented in the previous section. The unintegrated and integrated vertices for a particular nonconstant R-R field-strength will be computed, that are expected to be related to a nonconstant deformation of ten-dimensional superspace. At the end of the section, other two applications will be briefly described. No details will be given for these two, since they are not related to noncommutative geometry, which is the main subject of this thesis. However, the interested reader can refer to my work , where the three applications are discussed. #### 3.3.1 Motivation: Nonconstant superspace deformations In section 1.2.4, we have seen how non(anti)commutative superspaces, first studied in and in my paper , arise in the context of superstring theory, when open superstrings in the presence of R-R backgrounds and D-branes are considered. Up to now, only supergeometries with constant fermion-fermion anticommutators have been derived, associated to a constant R-R background, for the superstring compactified on a CY three-fold and for the uncompactified ten-dimensional superstring . In both cases the covariant formulation for the superstring has been used. In section 1.3 we have discussed how, in the bosonic case, nonconstant deformations of the coordinate algebra are related to the presence of a general curved NS-NS background. Therefore, it is natural to expect that more general, nonconstant R-R backgrounds can lead to nonconstant superspace deformations. In particular, in it was conjectured that from non-constant RR field strengths one can derive new equal-time commutation relations between coordinates $`x^m`$ and $`\theta ^\alpha `$ living on the boundaries such as $$\{\theta ^\alpha ,\theta ^\beta \}=\gamma _m^{\alpha \beta }x^m,$$ (3.200) generalizing the construction of Lie-algebraic non-commutative geometries to supermanifolds (for a different example of a Lie-algebraic geometry in superspace see ). The vertex operator for non-constant R-R fields strengths is the basic ingredient of this kind of analysis. #### 3.3.2 The ansatz for the RR field strength Applying the iterative method I introduced in the previous section, I will show how to compute the vertex for linearly $`x`$-dependent RR field strengths. This is the most simple $`x`$-dependent background one can consider. Moreover, it is interesting since it is supposed to be related to the nonconstant deformation (3.200). We will consider the following ansatz for the R-R field strength $$P^{\alpha \widehat{\beta }}=f^{\alpha \widehat{\beta }}+𝒞_m^{\alpha \widehat{\beta }}x^m$$ (3.201) where $`𝒞_m^{\alpha \widehat{\beta }}`$ is constant. $`P^{\alpha \widehat{\beta }}`$ must satisfy equations (3.108), which become $`\gamma _{\alpha \beta }^m𝒞_m^{\beta \widehat{\gamma }}=\gamma _{\widehat{\alpha }\widehat{\beta }}^m𝒞^{\gamma \widehat{\alpha }}=0`$ for the specific ansatz (3.201). Equations (3.108) can be rewritten in terms of forms by decomposing $`P^{\alpha \widehat{\beta }}`$ according to Dirac equations. For example, for type IIB we have the 1-form $`P_m`$, the 3-form $`P_{[mnp]}`$ and the 5-form $`F_{[mnpqr]}`$. Solving the Bianchi identities we get $`P_m=_mA`$, $`P_{[mnp]}=_{[m}A_{np]}`$,… and the field equations are $`^mP_m=^2A=0`$, $`^mP_{[mnp]}=^m_{[m}A_{np]}`$,… These can be solved in terms of quadratic polynomials $`A(x)=(10a_{(mn)}a_r^r\eta _{mn})x^mx^n`$, $`A_{[mn]}=(10a_{[mn],(rs)}a_{[mn,t]}^t\eta _{rs})x^rx^s`$,… where $`a_{(mn)}`$, $`a_{[mn],(rs)}`$,… are constant background fields. In the constant field strength case, our iterative procedure can be applied to compute the integrated and unintegrated vertex operators. One obtains $$𝒱_{z\overline{z}}^{(0,0)}=q_\alpha f^{\alpha \widehat{\beta }}\widehat{q}_{\widehat{\beta }},$$ (3.202) where $`q_\alpha `$ and $`\widehat{q}_{\widehat{\beta }}`$ are the supersymmetry currents given in (3.39). So it is easy to see that equation (3.56) is verified with $$𝒱^{(1,1)}=\chi _\alpha f^{\alpha \widehat{\beta }}\widehat{\chi }_{\widehat{\beta }},$$ (3.203) which is clearly BRST invariant (see (3.41) and (3.42)). Since in the $`\theta `$ and $`\widehat{\theta }`$ expansions of the physical and auxiliary superfields $`A_{\alpha \widehat{\beta }}`$,…,$`P^{\alpha \widehat{\beta }}`$ (see eqs. (3.154 and (3.178)) the number of bosonic derivatives acting on physical zero-order components grows with growing order in $`\theta `$ and $`\widehat{\theta }`$, it is clear that the ansatz (3.201) will correspond to only a few non-zero terms in the expansion. Actually, the highest-order contributions are $`\theta ^4\widehat{\theta }^2`$ and $`\theta ^2\widehat{\theta }^4`$ terms. Here we give the explicit expressions $`A_{\alpha \widehat{\beta }}`$ $`=`$ $`{\displaystyle \frac{1}{9}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\beta (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}(f^{\beta \widehat{\gamma }}+𝒞_n^{\beta \widehat{\gamma }}x^n)`$ (3.204) $`+`$ $`{\displaystyle \frac{1}{180}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\gamma }}(\gamma _q\widehat{\theta })_{\widehat{\delta }}𝒞_r^{\gamma \widehat{\delta }}`$ (3.205) $`+`$ $`{\displaystyle \frac{1}{180}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\delta (\gamma ^{nr}\theta )^\delta (\gamma _n\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^p\widehat{\theta })_{\widehat{\gamma }}𝒞_r^{\gamma \widehat{\gamma }},`$ (3.206) $`A_{\alpha p}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\beta (\gamma _p\widehat{\theta })_{\widehat{\gamma }}(f^{\beta \widehat{\gamma }}+𝒞_n^{\beta \widehat{\gamma }}x^n)`$ (3.207) $`+`$ $`{\displaystyle \frac{1}{36}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\beta }}(\gamma _q\widehat{\theta })_{\widehat{\gamma }}𝒞_r^{\gamma \widehat{\gamma }}`$ (3.208) $`+`$ $`{\displaystyle \frac{1}{60}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\beta (\gamma ^{nr}\theta )^\beta (\gamma _n\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}𝒞_r^{\gamma \widehat{\beta }}`$ (3.209) $`A_{m\widehat{\beta }}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(\gamma _m\theta )_\beta (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}(f^{\beta \widehat{\gamma }}+𝒞_n^{\beta \widehat{\gamma }}x^n)`$ (3.210) $`+`$ $`{\displaystyle \frac{1}{36}}(\gamma _m\theta )_\alpha (\gamma ^{nr}\theta )^\alpha (\gamma _n\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^p\widehat{\theta })_{\widehat{\gamma }}𝒞_r^{\gamma \widehat{\gamma }}`$ (3.211) $`+`$ $`{\displaystyle \frac{1}{60}}(\gamma _m\theta )_\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}(\gamma ^{rs}\widehat{\theta })^{\widehat{\gamma }}(\gamma _r\widehat{\theta })_{\widehat{\delta }}𝒞_s^{\gamma \widehat{\delta }}`$ (3.212) $`E_\alpha ^{\widehat{\beta }}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma (f^{\gamma \widehat{\beta }}+𝒞_n^{\gamma \widehat{\beta }}x^n)`$ (3.213) $`+`$ $`{\displaystyle \frac{1}{12}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}𝒞_q^{\gamma \widehat{\gamma }}`$ (3.214) $`+`$ $`{\displaystyle \frac{1}{60}}(\gamma ^m\theta )_\alpha (\gamma _m\theta )_\beta (\gamma ^{nr}\theta )^\beta (\gamma _n\theta )_\gamma 𝒞_r^{\gamma \widehat{\beta }}`$ (3.215) $`E_{\widehat{\beta }}^\alpha `$ $`=`$ $`{\displaystyle \frac{1}{3}}(\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}(f^{\alpha \widehat{\gamma }}+𝒞_m^{\alpha \widehat{\gamma }}x^m)`$ (3.216) $`+`$ $`{\displaystyle \frac{1}{12}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma (\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}𝒞_n^{\gamma \widehat{\gamma }}`$ (3.217) $`+`$ $`{\displaystyle \frac{1}{60}}(\gamma ^p\widehat{\theta })_{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\gamma }}(\gamma _q\widehat{\theta })_{\widehat{\delta }}𝒞_r^{\alpha \widehat{\delta }}`$ (3.218) $`A_{mp}`$ $`=`$ $`(\gamma _m\theta )_\beta (\gamma _p\widehat{\theta })_{\widehat{\gamma }}(f^{\beta \widehat{\gamma }}+𝒞_n^{\beta \widehat{\gamma }}x^n)`$ (3.219) $`+`$ $`{\displaystyle \frac{1}{12}}(\gamma _m\theta )_\beta (\gamma ^{nr}\theta )^\beta (\gamma _n\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}𝒞_r^{\gamma \widehat{\beta }}`$ (3.220) $`+`$ $`{\displaystyle \frac{1}{12}}(\gamma _m\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^{rs}\widehat{\theta })^{\widehat{\beta }}(\gamma _r\widehat{\theta })_{\widehat{\gamma }}𝒞_s^{\gamma \widehat{\gamma }}`$ (3.221) $`E_m^{\widehat{\beta }}`$ $`=`$ $`(\gamma _m\theta )_\gamma (f^{\gamma \widehat{\beta }}+𝒞_n^{\gamma \widehat{\beta }}x^n)`$ (3.222) $`+`$ $`{\displaystyle \frac{1}{4}}(\gamma _m\theta )_\gamma (\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}𝒞_q^{\gamma \widehat{\gamma }}`$ (3.223) $`+`$ $`{\displaystyle \frac{1}{12}}(\gamma _m\theta )_\alpha (\gamma ^{nr}\tau )^\alpha (\gamma _n\theta )_\gamma 𝒞_r^{\gamma \widehat{\beta }}`$ (3.224) $`E_p^\alpha `$ $`=`$ $`(\gamma _p\widehat{\theta })_{\widehat{\gamma }}(f^{\alpha \widehat{\gamma }}+𝒞_m^{\alpha \widehat{\gamma }}x^m)`$ (3.225) $`+`$ $`{\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma (\gamma _p\widehat{\theta })_{\widehat{\beta }}𝒞_n^{\gamma \widehat{\beta }}`$ (3.226) $`+`$ $`{\displaystyle \frac{1}{12}}(\gamma _p\widehat{\theta })_{\widehat{\beta }}(\gamma ^{qr}\widehat{\theta })^{\widehat{\beta }}(\gamma _q\widehat{\theta })_{\widehat{\gamma }}𝒞_r^{\alpha \widehat{\gamma }}`$ (3.227) $`P^{\alpha \widehat{\beta }}`$ $`=`$ $`(f^{\alpha \widehat{\beta }}+𝒞_m^{\alpha \widehat{\beta }}x^m)`$ (3.228) $`+`$ $`{\displaystyle \frac{1}{4}}(\gamma ^{mn}\theta )^\alpha (\gamma _m\theta )_\gamma 𝒞_n^{\gamma \widehat{\beta }}`$ (3.229) $`+`$ $`{\displaystyle \frac{1}{4}}(\gamma ^{pq}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}𝒞_q^{\alpha \widehat{\gamma }}.`$ (3.230) #### 3.3.3 The vertex for linearly $`x`$-dependent RR field strength To obtain the vertices $`𝒱^{(1,1)}`$ and $`𝒱_{z\overline{z}}^{(0,0)}`$ for the linearly $`x`$-dependent RR field strength we have to insert (3.204) and (3.230) back into (3.83) and (3.89). For the unintegrated vertex operator we find $`𝒱^{(1,1)}=\chi _\alpha f^{\alpha \widehat{\beta }}\widehat{\chi }_{\widehat{\beta }}`$ (3.231) $`+\chi _\alpha \left[\left(x^m\delta _\gamma ^\alpha \delta _{\widehat{\gamma }}^{\widehat{\beta }}+{\displaystyle \frac{1}{20}}(\gamma ^{qm}\widehat{\theta })^{\widehat{\beta }}(\gamma _q\widehat{\theta })_{\widehat{\gamma }}\delta _\gamma ^\alpha +{\displaystyle \frac{1}{20}}(\gamma ^{nm}\theta )^\alpha (\gamma _n\theta )_\gamma \delta _{\widehat{\gamma }}^{\widehat{\beta }}\right)𝒞_m^{\gamma \widehat{\gamma }}\right]\widehat{\chi }_{\widehat{\beta }}`$ (3.232) while for the integrated vertex operator $`𝒱_{z\overline{z}}^{(0,0)}`$ we obtain $`𝒱_{z\overline{z}}^{(0,0)}`$ $`=`$ $`q_\alpha f^{\alpha \widehat{\beta }}q_{\widehat{\beta }}`$ (3.234) $`+`$ $`q_\alpha \left[x^s\delta _\gamma ^\alpha \delta _{\widehat{\gamma }}^{\widehat{\beta }}+{\displaystyle \frac{1}{4}}(\gamma ^{rs}\theta )^\alpha (\gamma _r\theta )_\gamma \delta _{\widehat{\gamma }}^{\widehat{\beta }}+{\displaystyle \frac{1}{4}}(\gamma ^{ps}\widehat{\theta })^{\widehat{\beta }}(\gamma _p\widehat{\theta })_{\widehat{\gamma }}\delta _\gamma ^\alpha \right]𝒞_s^{\gamma \widehat{\gamma }}\widehat{q}_{\widehat{\beta }}`$ (3.235) $`+`$ $`\left[{\displaystyle \frac{1}{6}}(x^m+{\displaystyle \frac{1}{10}}\theta \gamma ^m\theta )(\theta \gamma _m\gamma ^{rs}\theta )N^{rs}\right](\gamma _r\theta )_\alpha 𝒞_s^{\alpha \widehat{\beta }}\widehat{q}_{\widehat{\beta }}`$ (3.236) $`+`$ $`q_\alpha 𝒞_s^{\alpha \widehat{\beta }}(\gamma _r\widehat{\theta })_{\widehat{\beta }}\left[{\displaystyle \frac{1}{6}}(\overline{}x^p+{\displaystyle \frac{1}{10}}\widehat{\theta }\gamma ^p\overline{}\widehat{\theta })(\widehat{\theta }\gamma _p\gamma ^{rs}\widehat{\theta })\widehat{N}^{rs}\right]`$ (3.237) Unfortunately, the complicated structure of $`𝒱_{z\overline{z}}^{(0,0)}`$ prevents from a simple analysis of superspace deformations as in . In the future it would be nice to find a way to study superspace deformations deriving from nonconstant R-R backgrounds. The computation of this vertex is a first step, but it is clear that plenty of work is needed to understand how to use it to compute the way it deforms the supergeometry. ### 3.4 Other applications #### 3.4.1 Vertex operators with R-R gauge potentials In the presence of D-branes, one can ask which states couple to them and which vertex operators describe such interaction. As it was discussed in in the framework of RNS formalism, one has to construct the vertex operators for R-R fields in the asymmetric picture. In addition, a propagating closed string (i.e. with non vanishing momentum) emitted from a disk or a D-brane, has to be off-shell. Therefore, one needs to break the BRST invariance by allowing a non vanishing commutator with $`Q_{L,0}+Q_{R,0}`$ where $`Q_{L/R,0}`$ are the picture conserving parts of BRST charges in the RNS formalism. In particular in the authors construct a solution of $`[Q_{L,1}+Q_{R,1},W]=0`$, where $`W`$ is the vertex operator in the asymmetric picture. The off-shell vertex operators directly couple the R-R potentials to the worldvolume of the D-brane. In , we constructed analogous vertices for closed superstrings which do not satisfy the classical supergravity equations of motion, but modified superfield constraints. They allow for a description of the R-R gauge potentials, in contradistinction to the on-shell formalism case, where only the field strengths $`P^{\alpha \widehat{\beta }}`$ appear. First of all, there are some important differences. The two BRST charges $`Q_L`$ and $`Q_R`$ contain a single term and therefore the decomposition used in is not viable. In addition, there are no different pictures (in the usual sense) for a given vertex since there are no superghosts associated to local worldsheet supersymmetry. There is, however, the possibility of constructing two operators which resemble the picture lowering and raising operator , as we briefly discussed in section 3.1.4, but the implications of this new idea in the present context have not been explored yet. Nevertheless we can construct an off-shell formalism by considering the following combination of vertices with different ghost numbers: $$𝒱^{(2)}=𝒱^{(2,0)}+𝒱^{(1,1)}+𝒱^{(0,2)},$$ (3.238) where the notation $`𝒱^{(a,b)}`$ stands for vertex operators with the left ghost number $`a`$ and with the right ghost number $`b`$. The ghost number of the l.h.s. is just the sum of the ghost numbers. By expanding this in terms of the pure spinor ghost and by applying the modified condition $$[Q_L+Q_R,𝒱_2]=0$$ (3.239) instead of the usual two conditions for the left and right sectors, we showed that this leads to equations of motion that are deformations of the usual supergravity constraints. The way constraints are relaxed to go off-shell follows very closely the case of $`N=1`$ super-Yang-Mills presented in . By following a procedure analogous to the one described in section 3.2.4, a suitable gauge-fixing can be applied to eliminate auxiliary fields. The superfields can then be expanded in $`\theta `$ and $`\widehat{\theta }`$ and one finds that R-R gauge-potentials explicitly appear. The construction of vertices with R-R potentials in covariant formulation has been also discussed in . There the authors considered only the constant case. #### 3.4.2 Antifields and the kinetic terms for closed superstring field theory The linearized form of supergravity equations written in terms of the BRST charges of the pure spinor sigma model gives us the framework to analyze some aspects of closed string field theory action. As it is well-known, the action for closed string field theory has to take into account the presence of selfdual forms (for example the five form in type IIB supergravity). This can be done either by breaking explicitly the Lorentz invariance, or by admitting an infinite number of fields in the action . In we showed that this action can be indeed constructed by mimicking the bosonic closed string field theory action discussed in (and in references therein). To this purpose, we derived the set of antifields for the massless sector of closed string theory, we discussed the coupling of the fields to the antifields for a closed string field theory action and we finally proposed a kinetic term which leads to the correct equations of motion. We showed that we could easily account for new fields which nevertheless do not propagate and we checked that the action had the correct symmetries leading to the complete BV action for type IIA/IIB supergravity. Since in we only dealt with linearized supergravity equations, we did not discuss generalizations of Witten string field $``$-product for the open superstring. Similarly, it was outside the scope of to construct a full-fledged closed string field theory. ## Chapter 4 Conclusions and outlook In this thesis I have presented my papers , where I investigated aspects of noncommutative geometry and superstring theory. When I started my research activity, it was already known that field theories defined on a noncommutative space arise as a low energy description of D-brane dynamics in the presence of a constant NS-NS background. This had been proven for the bosonic string, the RNS superstring and the $`N=2`$ string. Moreover, it had been shown that extending this discussion to superstring theory in GS formalism, where spacetime fermions are present, the (anti)commutation relations involving the fermions are not modified by the constant NS-NS background. These string theory results had already induced a growing interest in noncommutative field theory and many results had been obtained. For instance, it was already known that noncommutative field theories with time-space noncommutativity display awkward features, such as acausality and nonunitarity, and that these ill-defined theories cannot be obtained as a low energy limit of string theory. Moreover, the behavior of noncommutative field theory with respect to Poincaré symmetry was well-known. Moyal noncommutative deformation breaks Lorentz-invariance but preserves translation invariance. The generalization of the string results concerning D-branes in a constant NS-NS background had been generalized to nonconstant backgrounds afterwards, to show that a Kontsevich-like product replaces Moyal product in this case. The deformation is associative when the background is a closed two-form and nonassociative otherwise. All these results are reviewed in the first chapter of my thesis, so that my work can be put into context. In modern physics symmetries have such an importance that a theory is actually defined in terms of its symmetries and its field content. If a classical theory constructed out of the chosen fields does not contain all the interactions that are allowed by its symmetry structure, the missing terms will be generated at the quantum level when the theory is renormalized. A unifying aspect of my contributions to the field of noncommutative geometry is the way symmetries can be implemented in noncommutative generalizations of known ordinary theories. Since the noncommutative generalization of a given theory is not unique, preserving its symmetries (or at least some of them) in its noncommutative generalization is the first criterium one should take into account in the evaluation of the various possibilities, given the importance symmetries have in the physics of our world. In my paper , written in collaboration with D. Klemm and S. Penati, I have considered the very special symmetry, known as supersymmetry, relating bosonic and fermionic degrees of freedom in a theory. This symmetry has not been observed in nature yet, however there are hopes that it will make itself manifest at higher energies. For people who believe that the ultimate, fundamental theory of reality is string theory, supersymmetry is expected to be one of the characteristics of nature at sufficiently high energy, since string theory requires it for consistency. Supersymmetric theories are better studied in a formalism known as superspace. Superspace is an extension of spacetime where bosonic coordinates are accompanied by fermionic ones. There supersymmetry is realized in the form of generalized translations. The geometry of superspace is not flat, since a nonvanishing torsion is present. In we studied the way supersymmetric theories can be deformed by implementing a non(anti)commutative geometry, without supersymmetry to be lost. We first considered four-dimensional superspace with a Minkowski metric. We wrote down the most general algebra for bosonic and fermionic coordinates of superspace and then required covariance with respect to supertraslations and associativity. We also required the reality properties of the spinors in (anti)commutative superspace to be still valid in the deformed superspace. This was the crucial difference with respect to the previous work . We found that nontrivial coordinate algebras are allowed that involve fermionic coordinates. We also noted that, for consistency with respect to supersymmetry, turning on the anticommutators between the fermions implies that also terms depending on the fermionic coordinates appear in the fermion-boson and boson-boson commutators. The “trivial” superspace deformation where only bosonic coordinates are rendered noncommutative with a constant commutator is found as a special case. However, in Minkowski signature deformations involving nonzero fermion-fermion anticommutators are ruled out because of spinor reality conditions together with associativity requirements. Therefore, it is clear that more general deformations are allowed if spinor conjugation relations can be relaxed. In we noted that this is possible when moving to a four-dimensional superspace with Euclidean signature, that can only be defined when extended supersymmetry is present. By applying our procedure to Euclidean $`N=2`$ superspace, we found that deformations with constant fermion-fermion anticommutators are allowed. In we also discussed how to construct a non(anti)commutative $``$ product between superfields. Since our supergeometries involve coordinate-dependent terms in the coordinate algebra, a Moyal-like product is not associative, because of the superspace nontrivial torsion. However, we suggested that a Kontsevich-like product can be constructed, with the property of being associative if and only if the supercoordinate algebra is also associative. We gave the formula for the product up to second order in the deformation parameter and argued that there was no objection of principle to extending Kontsevich procedure to all orders. This was done later in . In we also noted that when superspace is rendered non(anti)commutative the supersymmetry algebra is deformed by terms quadratic in bosonic derivatives. These terms do not affect the coordinate algebra, but modify the supersymmetry transformation of a general superfield. Finally, we made some speculations on how the deformed superspaces we found could emerge from string theory, formulated in a manifestly target-space supersymmetric way (i.e. in Green-Schwarz or Berkovits formalisms ). Indeed, it was found that deformed superspaces similar to the ones I studied in naturally arise when the open superstring in the manifestly superPoincaré covariant formalism introduced by Berkovits is compactified on a CY three-fold in the presence of D-branes and a constant R-R background . The nonanticommutative superspace found there is related to the one I studied in the $`N=2`$ euclidean case by a change of variables and a reduction to $`N=1`$ by identification of the two fermionic coordinates on the boundary of the string worldsheet. Since in this superspace only one of the two Weyl spinor supersymmetry charges are deformed by quadratic terms in bosonic derivatives, Seiberg called it $`N=\frac{1}{2}`$ superspace. The constant R-R background considered in is allowed only in an euclidean signature. This is the stringy counterpart of the algebraic discussion in my paper , justifying by a geometrical argument why superspace deformations involving nonzero fermion-fermion anticommutators can only appear in an Euclidean signature. Deformed superspaces were also found to emerge in the uncompactified ten dimensional superstring when a constant R-R background and D-branes are present . However, since this background is not an exact solutions to the string equations, but only to the linearized equations, it is not obvious that this deformation survives the zero-slope limit necessary to obtain the low energy D-brane dynamics. After the discovery that non(anti)commutative superspaces can be obtained from the superstring, a lot of efforts in the study of non(anti)commutative field theories have been done and many interesting properties of $`N=\frac{1}{2}`$ Wess-Zumino model and super-Yang-Mills theory have been elucidated. Different deformations of theories with extended supersymmetry have been considered. The connection between deformed superspaces and supersymmetric matrix models has been elucidated. Up to now only constant R-R backgrounds have been considered. It would be interesting to study what happens in the more general, nonconstant case. A first step in this direction has been taken in my paper , written in collaboration with P.A. Grassi, where superstring vertex operators for linearly $`x`$-dependent R-R field strength have been computed. These objects are the main ingredient to generalize the string analysis presented in to the nonconstant case. Linearly $`x`$-dependent R-R backgrounds are expected to be related to a special kind of superspace deformation, with a Lie algebraic structure where the anticommutator between two fermions gives the bosonic coordinate. This can be interpreted by saying that spacetime bosonic coordinates have a fermionic substructure. In these terms, this issue was investigated in . Now I would like to go back to the main theme unifying my work, the implementation and implications of symmetries in noncommutative theories. Up to know I have discussed how I approached the problem of deformation of supersymmetric theories and the developments that followed my work, in both field and string theory. In two papers of mine , I considered instead the problem of constructing the noncommutative generalization of field theories possessing an infinite number of conserved local symmetry charges, i.e classically integrable. As it is well-known, in the commutative case the underlying symmetry structure of these theories has strong implications on their dynamics. In particular, two dimensional integrable theories have a factorized S-matrix and particle production does not occur. Most ordinary integrable theories in two and three dimensions have been shown to be related to a four dimensional theory, selfdual Yang-Mills, from which they can be obtained by dimensional reduction. The S-matrix of selfdual Yang-Mills also displays a peculiar property, all tree-level amplitudes beyond three-point being vanishing. This property has been used as a definition of integrability in four dimensions. Selfdual Yang-Mills is related to the $`N=2`$ string, characterized by an $`N=2`$ worldsheet supersymmetry. It has been proven that tree-level $`N=2`$ string dynamics coincides with selfdual Yang-Mills theory. The $`N=2`$ string can be coupled to a constant NS-NS background. In the presence of D-branes it has been shown that the low energy limit of the brane dynamics can be described by noncommutative selfdual Yang-Mills theory. Therefore, given an ordinary integrable two-dimensional field theory, one can first try to construct a noncommutative deformation that preserves classical integrability, in the sense of having an infinite number of conserved local charges (local in the sense that they are not written in terms of integrals). Then one can check if the usual S-matrix properties are still valid after the deformation. This is not obvious, since noncommutativity introduces nonlocality in the theory and, in the two-dimensional case, this necessarily affects the time coordinate, causing in general an acausal behavior and the breakdown of unitarity. Furthermore, one can investigate whether the theory can be obtained by reduction from noncommutative selfdual Yang-Mills, or even use this method to construct the two-dimensional theory in the first place. Another interesting issue is the study of solitons solutions. Commutative integrable theories usually display this kind of solutions. Noncommutative theories have new soliton solutions disappearing in the commutative limit, that exist thanks to the nonlocality introduced by Moyal product. Therefore, noncommutative integrable theories are an interesting setting to study both kinds of solitons. Noncommutative solitons are a field theory realization of the D-branes that are present in the string theory from which they have emerged in the low energy limit, thus the study of noncommutative solitons can give clues on D-brane dynamics. In , all these issues have been studied in the special case of the sine-Gordon model. This is an integrable two-dimensional field theory, describing the dynamics of a scalar field selfinteracting through an oscillating potential. Apart from the many general nice properties related to its integrability, the sine-Gordon model also exhibits very nice renormalization properties. Moreover, the ordinary sine-Gordon model is related to the Thirring model through bosonization. The relation between these two theories is a simple example of duality relating the weak coupling limit of a theory with the strong-coupling limit of the other. Bosonization has been studied in noncommutative geometry and it has been shown that the ordinary abelian $`U(1)`$ case is modified in the noncommutative setting so that a free fermion is not related to a free scalar, but to a scalar governed by a $`U(1)`$ WZW model. In , M.T. Grisaru and S. Penati have shown that an integrable noncommutative version of the sine-Gordon model can be constructed, where two equations govern the dynamics of a scalar (and generically complex) field. One equation reduces to the sine-Gordon equation in the commutative limit, the other vanishes in the commutative limit, has the form of a conservation law and in fact gives the first of the infinite conserved currents. This system does not seem to be overconstrained, since the class of its localized solutions is at least as large as the one in ordinary sine-Gordon. The doubling of the equations of motion is related to the fact that the $`U(1)`$ factor in the noncommutative group $`U(n)`$ does not decouple, in contradistinction to the ordinary case. Since commutative sine-Gordon theory is obtained from zero curvature conditions for gauge connections in $`SU(2)`$ group, in the noncommutative case this has to be enlarged to $`U(2)`$, which causes an additional equation of motion to appear. In my paper , written in collaboration with M.T. Grisaru, L. Mazzanti and S. Penati, we showed that the equations describing this noncommutative version of the sine-Gordon model can be obtained by dimensional reduction from noncommutative selfdual Yang-Mills in Yang formulation. Unfortunately, an action cannot be obtained by an analogous reduction procedure. However, in , we found an action that gives the correct equations of motion. We then studied tree-level scattering amplitudes and found that acausality is present and the S-matrix is not factorized. Therefore, it seems that in general the presence of an infinite number of conserved currents is not sufficient for the S-matrix to be factorized in noncommutative case. We finally discussed the relation of our model with the noncommutative $`U(1)`$ Thirring model. The fact that dimensional reduction does not work at the level of action, but only at the level of the equations of motion, is a sign that the parametrization of the degrees of freedom we considered in is not the correct one. In , in collaboration with O. Lechtenfeld, L. Mazzanti, S. Penati and A.D. Popov, we considered a different noncommutative generalization of the sine-Gordon model, that is also obtained by dimensional reduction from noncommutative selfdual Yang-Mills, by considering the $`U(1)\times U(1)`$ subgroup of $`U(2)`$. While in our first attempt to construct a noncommutative integrable sine-Gordon model we had modified the kinetic term of the ordinary theory to a $`U(1)`$ WZW-like term, while maintaining the usual cosine structure of the interaction term, in this new noncommutative theory we modified also the interaction term structure, in a way that two real scalars parametrizing the $`U(1)\times U(1)`$ subgroup of $`U(2)`$ are coupled by a nontrivial term. Again, the theory is described by two equations of motion. One of the two becomes trivial in the commutative limit, while the other one gives the ordinary sine-Gordon equation. Solitons solutions of this model have been studied by making use of the dressing method. Moreover, tree-level scattering amplitudes have been computed and proven to be factorized and causal, in spite of the presence of time-space noncommutativity. Therefore this second noncommutative version of the sine-Gordon model has inherited the nice classical properties of its ordinary counterpart. It would be interesting to move on to a quantum description of the model to study its renormalizability properties and to investigate whether integrability survives at the quantum level. A quite striking feature of the noncommutative sine-Gordon system we constructed is that its tree-level amplitudes are completely independent of the noncommutativity parameter. It would be interesting to understand why this happens, for instance if this is a general feature of noncommutative integrable systems or a peculiar property of the sine-Gordon model. Another aspect that still needs to be investigated is the relation of this $`U(1)\times U(1)`$ theory with noncommutative fermions models. Somehow it would be more natural to see what happens in the bosonization of $`U(2)`$ fermion models and then consider a reduction to $`U(1)\times U(1)`$. Another possible development would be the study of other kinds of duality in this system, for instance T-duality. As I anticipated when discussing the superstring theory origin of the deformed superspaces I introduced in , the covariant formulation for the superstring has been proven to be superior in dealing with many string issues that were not treatable with the other known formalisms. In particular this was true for deformed superspaces, that were shown to arise in the presence of R-R backgrounds. These cannot be dealt with when using the RNS formalism, while they can be treated with the GS formalism, that however is quite clumsy because of the lack of manifest Lorentz covariance. Because of this connection with the work I did in the field of noncommutative geometry, I started to study this formalism. As I previously said, to generalize the results of to a more general, nonconstant background one needs to determine the corresponding vertex operator first. This is not an easy task in the covariant formulation, since the great amount of manifest symmetry makes the formalism redundant. The vertices are written in terms of superfields that have to satisfy the linearized supergravity equations of motion. In my paper , written in collaboration with P.A. Grassi, we described an iterative procedure to compute the vertex operators in terms of the physical fields only. To do this a suitable gauge fixing must be imposed that removes the auxiliary fields from the vertices. We used this technique to compute the vertices for linearly $`x`$-dependent R-R field strength that may be related to a Lie-algebraic deformation of superspace, as I already said, but we also discussed other two applications of our analysis that are a little bit out of the path, having no direct connection to noncommutative geometry. We showed how an off-shell formulation of the superstring vertices can be constructed, how the corresponding equations of motion are a deformation of the usual superfield constraints and how the gauge fixing necessary to implement our iterative procedure can be applied to this case. The motivation for this discussion relies in the fact that a propagating closed string emitted from a D-brane has to be off-shell. The off-shell vertex operators couple the R-R potential (and not the field strength) to the worldvolume of the brane. In our discussion, we showed that the vertices we obtain explicitly contain the R-R potential. Therefore our analysis could be useful in the study of D-brane dynamics in the covariant formalism. Finally, we used our framework to determine the antifield equations of motion and to make a proposal for the kinetic term of closed superstring field theory. Since this kinetic term is written in superspace formulation, the first thing one should now do to make nontrivial checks about its properties is to rewrite it in component formulation. To conclude, in my work I have mostly investigated the way one can implement noncommutativity in theories with a special symmetry structure, such as supersymmetric and classically integrable theories. Requiring the symmetries to be preserved in the noncommutative deformation is a strong constraint that allows to make a selection between the many possibile noncommutative versions of the same theory. In the case of integrable theories, I explicitly discussed the case of the sine-Gordon theory and eventually found its noncommutative generalization that has infinite conserved currents and also displays all the nice features that in ordinary theories are implied by the symmetry structure, such as factorization of the S-matrix. However, many issues still have to be investigated, such as for instance the connection between the noncommutative sine-Gordon system and $`U(2)`$ fermion models and the reason for the complete absence of a dependence on the noncommutativity parameter in the tree-level amplitudes. In the case of supersymmetric theories, I constructed superspace deformations that preserve supersymmetry, are associative and respect the usual spinor reality properties. The connection between the deformations I found and superstring theory in the presence of R-R backgrounds, described in the covariant formalism, has led me to deepen my knowledge of this formulation of superstring theory. Even if this might seem constructed “ad hoc”, without a strong underlying principle, it is the first superPoincaré covariant formulation of the superstring that works and it has already proven to very suitable to handle R-R backgrounds and to prove general theorems about superstring amplitudes. Since the computation of superstring vertex operators is not an easy task in the covariant formalism, I provided a recursive technique to compute the vertices in terms of physical fields only. This analysis of vertex operators has many possible applications. I discussed applications to the computation of superspace deformations in the presence of nonconstant R-R backgrounds, to the computation of vertices in the off-shell formulation, that are useful in the study of D-brane dynamics, and to the construction of kinetic terms for a closed superstring field theory. Since the covariant superstring in the presence of NS-NS backgrounds has not been studied in detail yet, it would be very interesting to perform an analysis of the string origin of noncommutative geometry in this formalism. Indeed, since NS-NS and R-R backgrounds are treated in a very “symmetric” way in this formulation of the superstring, this setting should be very useful to study S-duality. ## Appendix A Conventions ### A.1 Superspace conventions in $`d=4`$ In four dimensions, $`N=1`$ Minkowski superspace is described by a set of coordinates $`Z^A(x^{\alpha \dot{\alpha }},\theta ^\alpha ,\overline{\theta }^{\dot{\alpha }})`$, where $`x^{\alpha \dot{\alpha }}x^{\underset{¯}{a}}`$ are the four bosonic real coordinates, while $`\theta ^\alpha `$ (and $`\overline{\theta }^{\dot{\alpha }}=(\theta ^\alpha )^{}`$) are complex two–component Weyl fermions. We use conventions of Superspace , with $`(\psi ^\alpha )^{}=\overline{\psi }^{\dot{\alpha }}`$, $`(\psi _\alpha )^{}=\overline{\psi }_{\dot{\alpha }}`$. The supersymmetry algebra $`\{Q_\alpha ,\overline{Q}_{\dot{\alpha }}\}=P_{\alpha \dot{\alpha }}`$ $`\{Q_\alpha ,Q_\beta \}=\{\overline{Q}_{\dot{\alpha }},\overline{Q}_{\dot{\beta }}\}=0`$ $`[P_{\underset{¯}{a}},P_{\underset{¯}{b}}]=0`$ (A.1) with $`\overline{Q}_{\dot{\alpha }}=Q_\alpha ^{}`$, is realized by $`Q_\alpha =i\left(_\alpha {\displaystyle \frac{i}{2}}\overline{\theta }^{\dot{\alpha }}_{\alpha \dot{\alpha }}\right)`$ $`\overline{Q}_{\dot{\alpha }}=i\left(\overline{}_{\dot{\alpha }}{\displaystyle \frac{i}{2}}\theta ^\alpha _{\alpha \dot{\alpha }}\right)`$ $`P_{\alpha \dot{\alpha }}=i_{\alpha \dot{\alpha }}`$ (A.2) Under supersymmetry transformations a generic superfield $`V`$ transforms according to $`\delta V=i\left(ϵ^\alpha Q_\alpha +\overline{ϵ}^{\dot{\alpha }}\overline{Q}_{\dot{\alpha }}\right)V`$. In particular, the action on the superspace coordinates defines supertranslations $`\delta x^{\underset{¯}{b}}i\left(ϵ^\alpha Q_\alpha +\overline{ϵ}^{\dot{\alpha }}\overline{Q}_{\dot{\alpha }}\right)x^{\underset{¯}{b}}={\displaystyle \frac{i}{2}}\left(ϵ^\beta \overline{\theta }^{\dot{\beta }}+\overline{ϵ}^{\dot{\beta }}\theta ^\beta \right)`$ $`\delta \theta ^\beta i\left(ϵ^\alpha Q_\alpha +\overline{ϵ}^{\dot{\alpha }}\overline{Q}_{\dot{\alpha }}\right)\theta ^\beta =ϵ^\beta `$ $`\delta \overline{\theta }^{\dot{\beta }}i\left(ϵ^\alpha Q_\alpha +\overline{ϵ}^{\dot{\alpha }}\overline{Q}_{\dot{\alpha }}\right)\overline{\theta }^{\dot{\beta }}=\overline{ϵ}^{\dot{\beta }}`$ (A.3) Covariant derivatives with respect to (A.3) are $`D_A(D_\alpha ,\overline{D}_{\dot{\alpha }},_{\alpha \dot{\alpha }})`$ where $`D_\alpha =_\alpha +{\displaystyle \frac{i}{2}}\overline{\theta }^{\dot{\alpha }}_{\alpha \dot{\alpha }}`$ , $`\overline{D}_{\dot{\alpha }}=\overline{}_{\dot{\alpha }}+{\displaystyle \frac{i}{2}}\theta ^\alpha _{\alpha \dot{\alpha }}`$ $`_{\alpha \dot{\alpha }}`$ $`=`$ $`i\{D_\alpha ,\overline{D}_{\dot{\alpha }}\}`$ (A.4) Moreover, they satisfy $`\{D_\alpha ,D_\beta \}=\{\overline{D}_{\dot{\alpha }},\overline{D}_{\dot{\beta }}\}=0`$ and anticommute with the generators of supersymmetry transformations. We can define left and right grassmannian derivatives according to the following rules $`(_L)_\alpha \theta ^\beta \stackrel{}{}_\alpha \theta ^\beta =\delta _\alpha ^\beta ,(\overline{}_L)_{\dot{\alpha }}\overline{\theta }^{\dot{\beta }}\stackrel{}{\overline{}}_{\dot{\alpha }}\overline{\theta }^{\dot{\beta }}=\delta _{\dot{\alpha }}^{\dot{\beta }}`$ $`(_R)_\alpha \theta ^\beta \theta ^\beta \stackrel{}{}_\alpha =\delta _\alpha ^\beta ,(\overline{}_R)_{\dot{\alpha }}\overline{\theta }^{\dot{\beta }}\overline{\theta }^{\dot{\beta }}\stackrel{}{\overline{}}_{\dot{\alpha }}=\delta _{\dot{\alpha }}^{\dot{\beta }}`$ (A.5) Their action on a generic superfield is defined as $$_LV\stackrel{}{}V,_RVV\stackrel{}{}$$ (A.6) Notice that these definitions hold independently of the nature of the superfield V. In particular, in the case of a spinorial superfield $`V_\beta `$ we have $`(_R)_AV_\beta V_\beta \stackrel{}{}_A`$. As a consequence of the general definitions (A.5, A.6) we immediately obtain $`_LV=_RV`$ for any tensorial superfield, whereas $`_LV_\beta =_RV_\beta `$ for any spinorial superfield. From the identities $`\left(\stackrel{}{}_\alpha \theta ^\beta \right)^{}=\overline{\theta }^{\dot{\beta }}\stackrel{}{\overline{}}_{\dot{\alpha }},\left(\stackrel{}{\overline{}}_{\dot{\alpha }}\overline{\theta }^{\dot{\beta }}\right)^{}=\theta ^\beta \stackrel{}{}_\alpha `$ $`\left(\stackrel{}{}^\alpha \theta ^\beta \right)^{}=\overline{\theta }^{\dot{\beta }}\stackrel{}{\overline{}}^{\dot{\alpha }},\left(\stackrel{}{\overline{}}^{\dot{\alpha }}\overline{\theta }^{\dot{\beta }}\right)^{}=\theta ^\beta \stackrel{}{}^\alpha `$ (A.7) the hermitian conjugation rules for left and right derivatives follow $$((_L)_\alpha )^{}=(\overline{}_R)_{\dot{\alpha }},(_L^\alpha )^{}=(\overline{}_R)^{\dot{\alpha }}$$ (A.8) We can also introduce left and right bosonic derivatives which are simply given by $$(_L)_{\alpha \dot{\alpha }}x^{\beta \dot{\beta }}_{\alpha \dot{\alpha }}x^{\beta \dot{\beta }}=\delta _\alpha ^\beta \delta _{\dot{\alpha }}^{\dot{\beta }},(_R)_{\alpha \dot{\alpha }}x^{\beta \dot{\beta }}x^{\beta \dot{\beta }}\stackrel{}{}_{\alpha \dot{\alpha }}=\delta _\alpha ^\beta \delta _{\dot{\alpha }}^{\dot{\beta }}$$ (A.9) Their action on a superfield $`V`$ is defined as $`_LV\stackrel{}{}V`$ and $`_RVV\stackrel{}{}`$. Therefore, from (A.9), it easily follows that $`(_L)_{\alpha \dot{\alpha }}V=(_R)_{\alpha \dot{\alpha }}V`$ for any superfield. As a consequence of the previous identities, left and right covariant derivatives can be equally defined. Left covariant derivatives act on a generic superfield from the left as $$(D_L)_\alpha V\stackrel{}{D}_\alpha V,(\overline{D}_L)_{\dot{\alpha }}V\stackrel{}{\overline{D}}_{\dot{\alpha }}V$$ (A.10) where $`D_\alpha `$ and $`\overline{D}_{\dot{\alpha }}`$ are explicitly given in (A.4). Right covariant derivatives are defined as acting from the right $`(D_R)_\alpha VV\stackrel{}{D}_\alpha =V\left(\stackrel{}{}_\alpha +{\displaystyle \frac{i}{2}}\stackrel{}{}_{\alpha \dot{\alpha }}\overline{\theta }^{\dot{\alpha }}\right)`$ $`(\overline{D}_R)_{\dot{\alpha }}VV\stackrel{}{\overline{D}}_{\dot{\alpha }}=V\left(\stackrel{}{\overline{D}}_{\dot{\alpha }}=\stackrel{}{\overline{}}_{\dot{\alpha }}+{\displaystyle \frac{i}{2}}\stackrel{}{}_{\alpha \dot{\alpha }}\theta ^\alpha \right)`$ (A.11) It is easy to check that $`D_LV=D_RV`$ on any tensorial superfield, whereas $`D_LV_\beta =D_RV_\beta `$. Moreover, left and right derivatives are related by hermitian conjugation $`((D_L)_\alpha )^{}=(\overline{D}_R)_{\dot{\alpha }}`$ and $`((D_L)^\alpha )^{}=(\overline{D}_R)^{\dot{\alpha }}`$. Defining the right momentum operator as $`(P_R)_{\alpha \dot{\alpha }}V=iV\stackrel{}{}_{\alpha \dot{\alpha }}`$, it is easy to show that the algebra of right derivatives is the standard one $`\{D_R,D_R\}=\{\overline{D}_R,\overline{D}_R\}=0`$ and $`\{(D_R)_\alpha ,(\overline{D}_R)_{\dot{\alpha }}\}=(P_R)_{\alpha \dot{\alpha }}`$. Moreover $`[D_\alpha ^R,D_\beta ^L]=0=[\overline{D}_{\dot{\alpha }}^R,\overline{D}_{\dot{\beta }}^L]`$ $`[D_\alpha ^R,\overline{D}_{\dot{\alpha }}^L]V=(P_R)_{\alpha \dot{\alpha }}V`$ $`[D_\alpha ^L,\overline{D}_{\dot{\alpha }}^R]V=(P_L)_{\alpha \dot{\alpha }}V`$ (A.12) for any tensor superfield. When the commutators are applied to a spinor $`V_\beta `$, a minus sign appears on the r.h.s. of the last two identities due to the anticommutation of the spinorial derivatives with $`V_\beta `$. Following the same procedure one can equally define left and right supersymmetry generators as $`(Q_L)_AV\stackrel{}{Q}_AV`$ and $`(Q_R)_\alpha VV\stackrel{}{Q}_\alpha =V\left[i\left(\stackrel{}{}_\alpha {\displaystyle \frac{i}{2}}\stackrel{}{}_{\alpha \dot{\alpha }}\overline{\theta }^{\dot{\alpha }}\right)\right]`$ $`(\overline{Q}_R)_{\dot{\alpha }}VV\stackrel{}{\overline{Q}}_{\dot{\alpha }}=V\left[i\left(\stackrel{}{\overline{}}_{\dot{\alpha }}{\displaystyle \frac{i}{2}}\stackrel{}{}_{\alpha \dot{\alpha }}\theta ^\alpha \right)\right]`$ (A.13) The algebra of right generators is again given by (A.1). The algebra of the commutators on a tensor superfield is $`[Q_\alpha ^R,Q_\beta ^L]=0=[\overline{Q}_{\dot{\alpha }}^R,\overline{Q}_{\dot{\beta }}^L]`$ $`[Q_\alpha ^R,\overline{Q}_{\dot{\alpha }}^L]V=(P_L)_{\alpha \dot{\alpha }}V`$ $`[Q_\alpha ^L,\overline{Q}_{\dot{\alpha }}^R]V=(P_R)_{\alpha \dot{\alpha }}V`$ (A.14) Instead, when the commutators act on spinor objects we get a change of sign on the r.h.s. of the last two equalities. In euclidean signature a reality condition on spinors is applicable only in the presence of extended supersymmetry. In the simplest case, $`N=2`$ euclidean superspace, the two–component Weyl spinors satisfy a symplectic Majorana condition $$(\theta _i^\alpha )^{}=\theta _\alpha ^iC^{ij}\theta _j^\beta C_{\beta \alpha },(\overline{\theta }^{\dot{\alpha },i})^{}=\overline{\theta }_{\dot{\alpha },i}\overline{\theta }^{\dot{\beta },j}C_{\dot{\beta }\dot{\alpha }}C_{ji}$$ (A.15) with $`C^{12}=C_{12}=i`$. The choice of covariant derivatives and supersymmetry charges is the obvious generalization of the $`N=1`$ Minkowski case. ### A.2 Superspace conventions in $`d=10`$ We describe Minkowski ten dimensional $`N=2`$ superspace by the coordinates $`(x^m,\theta ^\alpha ,\widehat{\theta }^{\widehat{\beta }})`$. Depending whether one is in type IIA or IIB case, the spinors $`\theta ^\alpha `$ and $`\widehat{\theta }^{\widehat{\beta }}`$ have opposite or same chirality. In ten dimensions with Minkowski signature one can use Dirac matrices $`\mathrm{\Gamma }^m=\{I(i\tau _2),\sigma ^\mu \tau _1,\chi \tau _1\}`$, where $`m=0,\mathrm{},9`$ and $`\mu =1,\mathrm{},8`$. $`\tau _i`$, $`i=1,2,3`$ are the Pauli matrices, $`\sigma ^\mu `$ are eight real symmetric $`16\times 16`$ off-diagonal Dirac matrices in $`d=8`$ with euclidean signature, while $`\chi `$ is the real $`16\times 16`$ diagonal chirality matrix in $`d=8`$. The chirality matrix in $`d=10`$ with Minkowski signature is then $`I\tau _3`$ and the charge conjugation matrix $`C`$, satisfying $`C\mathrm{\Gamma }^m=(\mathrm{\Gamma }^m)^TC`$, is numerically equal to $`\mathrm{\Gamma }^0`$. If one uses spinors $`\mathrm{\Psi }^T=(\alpha _L,\beta _R)`$ with spinor indices $`\alpha _L^\alpha `$ and $`\beta _{R,\dot{\beta }}`$, the index structure of the Dirac matrices and the charge conjugation matrix is $$\mathrm{\Gamma }^m=\left(\begin{array}{cc}0& (\sigma ^m)^{\alpha \dot{\beta }}\\ (\stackrel{~}{\sigma }^m)_{\dot{\beta }\gamma }& 0\end{array}\right),C=\left(\begin{array}{cc}0& c_\alpha ^{\dot{\beta }}\\ c_\gamma ^{\dot{\beta }}& 0\end{array}\right)$$ (A.16) where $`\sigma ^m=\{I,\sigma ^\mu ,\chi \}`$ and $`\stackrel{~}{\sigma }^m=\{I,\sigma ^\mu ,\chi \}`$. The matrices $`c_\alpha ^{\dot{\beta }}`$ and $`c_\gamma ^{\dot{\beta }}`$ are numerically equal to $`I_{16\times 16}`$ and $`I_{16\times 16}`$, respectively. The matrices $`\gamma ^m`$ used in the text are given by $`\gamma ^{m\dot{\alpha }\dot{\beta }}=c_\beta ^{\dot{\alpha }}(\sigma ^m)^{\beta \dot{\beta }}`$ and $`\gamma _{\alpha \beta }^m=c_\alpha ^{\dot{\beta }}(\stackrel{~}{\sigma }^m)_{\beta \dot{\beta }}`$. From now on and in the text dots are omitted. The spinors $`\alpha _L`$ and $`\beta _R`$ form inequivalent representations of $`SO(9,1)`$. Spin indices cannot be raised or lowered with the charge conjugation matrix, but $`\alpha _L^\alpha c_\alpha ^{\dot{\beta }}\beta _{R,\dot{\beta }}`$ is Lorentz invariant. One finds that the $`16\times 16`$ symmetric matrices $`\gamma _{\alpha \beta }^m`$ and $`\gamma ^{m\alpha \beta }`$ satisfy $$\gamma _{\alpha \beta }^m\gamma ^n\beta \gamma +\gamma _{\alpha \beta }^n\gamma ^{m\beta \gamma }=2\eta ^{mn}\delta _\alpha ^\gamma $$ (A.17) and $$\gamma _{m(\alpha \beta }\gamma _{\gamma )\delta }^m=0$$ (A.18) which makes Fierz rearrangements very easy. Our conventions for $`d=10`$ $`N=2`$ superspace covariant derivatives and supersymmetry charges are $`D_\alpha =_\alpha +{\displaystyle \frac{1}{2}}(\gamma ^m\theta )_\alpha _m,Q_\alpha =_\alpha {\displaystyle \frac{1}{2}}(\gamma ^m\theta )_\alpha _m,`$ (A.19) $`\widehat{D}_{\widehat{\alpha }}=_{\widehat{\alpha }}+{\displaystyle \frac{1}{2}}(\gamma ^m\widehat{\theta })_{\widehat{\alpha }}_m,\widehat{Q}_{\widehat{\alpha }}=_{\widehat{\alpha }}{\displaystyle \frac{1}{2}}(\gamma ^m\theta )_{\widehat{\alpha }}_m,`$ (A.20) which satisfy $`\{D_\alpha ,D_\beta \}=\gamma _{\alpha \beta }^m_m,\{\widehat{D}_{\widehat{\alpha }},\widehat{D}_{\widehat{\beta }}\}=\gamma _{\widehat{\alpha }\widehat{\beta }}^m_m,\{D_\alpha ,\widehat{D}_{\widehat{\beta }}\}=0`$ (A.21) $`\{D_\alpha ,Q_\beta \}=0,\{\widehat{D}_{\widehat{\alpha }},\widehat{Q}_{\widehat{\beta }}\}=0.`$ (A.22) #### Acknowledgements First of all I would like to thank Silvia Penati so much for her continuous encouragement, help and support in these many years. I would like to thank the Department of Nuclear and Theoretical Physics in Pavia and INFN, sezione di Pavia, for financial support. In particular, I am very grateful to Annalisa Marzuoli, Mauro Carfora and Sergio Ratti for their very kind help and encouragement. I am very grateful to the Physics Department G. Occhialini, University of Milano-Bicocca, for the very kind hospitality throughout these years and for providing funding that supported my participation into some schools and conferences. I would like to thank all the “baretto” friends in particular for the funniest lunch breaks ever. I would like to thank the C. N. Yang Institute for Theoretical Physics at Stony Brook, U.S.A., for the very kind hospitality during my stay within the International Doctorate Programme between the State University of New York at Stony Brook and the University of Pavia. In particular, I would like to thank Martin Roček for his kind help and support during my stay there, for the many useful discussions and especially for coming to Italy and participating into my thesis defense as a member of the evaluating committee. I would like to thank Pietro Antonio Grassi, Marc Grisaru, Dietmar Klemm, Olaf Lechtenfeld, Liuba Mazzanti, Silvia Penati and Alexander Popov for giving me the opportunity to collaborate with them and to learn so much from them. Very special thanks go to Dietmar Klemm and Marc Grisaru for their very strong encouragement, help and support. I am especially grateful to Ulf Lindström for kindly accepting to be the external referee for this thesis, for his very careful reading of it and for the many useful discussions, comments and suggestions. Finally, I would like to thank my many ex officemates, Claudio Dappiaggi, Valeria Gili, Liuba Mazzanti, Marcello Musso, Alberto Romagnoni and Andrea Sartirana for tolerating my bad temper, loud voice, invasive plants, my mess and very bad habit of locking one of them in particular inside the office. Now I’m alone in my office and missing all of you.
warning/0506/astro-ph0506334.html	ar5iv	text	# Cosmological perturbations in multiple-field inflation ## I Introduction It is well known that the cosmological inflation, an epoch of accelarated expansion, provides a causal mechanism for the generation and evolution of large-scale structure formation in the Universe Guth ; Lindebook ; Liddlebook . Inflation explains a number of puzzles of the Big-Bang theory, such as homogeneity, the isotropy of the Cosmic Microwave Background Radiation (CMBR) and the flatness of space-like sections. As an added bonus, it connects cosmology with high-energy physics, thus forming a cosmic laboratory where one may probe physics beyond the Standard Model. The very high accuracy CMBR data recently obtained by the WMAP satellite WMAP provide a new impetus to compare the predictions of cosmological inflation with observations and hopefully to discover new physics in the very high energy regime. At the simplest level, the inflationary scenario is implemented by assuming that the matter is described by a single scalar field, the inflaton, which is a special case of a perfect fluid Guth ; Single . As the early Universe undergoes inflation, quantum fluctuations of the scalar fields are generated which become classical after crossing the event horizon. During the decelaration phase they re-enter the horizon and seed the observed density perturbations. Cosmological perturbations in a single field inflation has been thoroughly investigated in the past Liddlebook ; Mukhanovetal . Despite initial successes, the single component inflation also has its drawbacks. It was realized quite early that, in its original form, the inflationary scenario suffers from what is called the graceful-exit problem Guth ; Exit . Linde Linde90 showed that in order to achieve sufficient inflation consistent with the observed density perturbations, before the Universe exits from the inflationary epoch, one requires at least two scalar fields without modifying Einstein gravity, and without sacrificing natural initial conditions. Of course, there are other motivations for incorporating multiple scalar fields contributing to the dynamics of inflation. When constructing models of inflation inspired by particle physics theories such as low energy effective supergravity derived from superstrings, one obtains many scalar fields (see Lythrep for a recent review). Thus it is necessary to have a general framework for handling cosmological perturbations in a situation where the matter sector consists of an arbitrary number of scalar fields. A method for treating density perturbations in multicomponent inflation was proposed in Tent , but see also Stewart1 ; Stewart2 ; Stewart3 . The study of cosmological perturbations was initiated by Lifshitz Lifs in 1946 when he analyzed hydrodynamical fluid perturbations in Einstein gravity. He assumed a particular gauge which is now known as the synchronous gauge. This gauge does not completely fix the gauge degrees of freedom and the spurious gauge modes have to be properly sorted out in order to obtain correct results. Later on, the zero-shear gauge was used by Harrison Harris and the comoving gauge by Nariai Nariai . However, it was the seminal paper by Bardeen Bardeen80 which helped put cosmological perturbations on a proper footing. He introduced a number of gauge-invariant variables to the linear order, in terms of which the perturbations became much simpler to analyze. Reviews of cosmological perturbations may be found in Mukhanovetal and Kodamasasaki . A different approach to cosmological perturbations was elaborated by Hwang and colleagues in a series of papers Hwang1 ; Hwang2 ; Hwang3 ; Hwang4 , following a suggestion by Bardeen Bardeen88 , that rather than imposing a particular gauge condition right from the beginning, it is often advantageous to express the perturbations without specifying any gauge. This then adds the flexibility of adopting different gauge conditions at a much later stage, depending upon the nature of each problem. Moreover, it becomes easy to relate results between various gauge-dependent and gauge-invariant techniques. This approach has been termed the gauge-ready method. In this paper we apply the gauge-ready method to analyze perturbations in cosmological inflation driven by multicomponent scalar fields. This paper is organized as follows. In Section II we set up the equations describing multicomponent scalar fields with a non-trivial field metric coupled non-minimally to Einstein gravity. We then introduce a set of basis vectors, using Gram-Schmidt orthonormalization, which enables us to disentangle multiple-field effects from single-fields ones. The background equations, metric perturbations and the perturbed order variables are presented in Section III. Here we also discuss briefly the issue of gauge transformations as applied to cosmological perturbations. In Section IV we introduce the gauge-ready approach to cosmological perturbations in the multiple-field inflation scenario. We present the perturbation equations in the gauge-ready form, as well as in terms of gauge-invariant variables derived from the gauge-ready equations. Slow-roll variables in the context of multicomponent inflation are presented in Section V. We proceed to apply canonical quantization to the density perturbations. The solutions to the equations for quantized perturbations are then derived to the first order in slow-roll. A brief discussion of vector and tensor perturbations is also presented. We conclude in Section VI. ## II Preliminaries ### II.1 The scalar fields In this Section, we explain our notation and set up the basic equations needed for our analysis. As our starting point, we consider Einstein gravity coupled to an arbitrary number of real scalar fields. We write the Lagrangean as $``$ $`=`$ $`\sqrt{g}\left({\displaystyle \frac{1}{2\kappa _0^2}}R{\displaystyle \frac{1}{2}}^\mu \mathit{\varphi }_\mu \mathit{\varphi }V(\mathit{\varphi })\right)`$ (1) $`=`$ $`\sqrt{g}\left({\displaystyle \frac{1}{2\kappa _0^2}}R{\displaystyle \frac{1}{2}}g^{\mu \nu }_\mu \mathit{\varphi }^T𝑮_\nu \mathit{\varphi }V(\mathit{\varphi })\right).`$ Here $`R`$ is the scalar curvature, $`\kappa _0^28\pi G`$, and we set $`c=1`$. For the scalar fields we use a vector notation, $`\mathit{\varphi }(\varphi ^a)`$, where the indices $`a,b,c,\mathrm{}=1,2,3,`$ $`\mathrm{},N`$ label the $`N`$–components in field space. Further, $`g\text{det}(g_{\mu \nu })`$, and $`\mu ,\nu ,\mathrm{}`$ denote the spacetime indices. For repeated indices, the summation convention applies. The second quantity within the parentheses of Eq. (1) represents a nonlinear sigma-model like non-minimal kinetic term. Such a kinetic term appears in various models of high-energy physics Lythrep . Also $`V(\mathit{\varphi })`$ is an arbitrary scalar potential. Following the authors of Tent , we can interpret the scalars $`\mathit{\varphi }`$ as coordinates $`(\varphi ^a)`$ on a real manifold $``$ induced with a symmetric Riemannian metric $`𝑮`$ having components $`G_{ab}`$ in the field space. The field metric is chosen to be positive-definite so that the Hamiltonian is bounded from below. The special case of minimally-coupled fields corresponds to the situation $`G_{ab}\delta _{ab}`$. From the components $`G_{ab}`$ we can define the connection coefficients $`\mathrm{\Gamma }_{bc}^a`$ in the usual manner, $$\mathrm{\Gamma }_{bc}^a=\frac{1}{2}G^{ad}\left(G_{bd,c}+G_{cd,b}G_{bc,d}\right).$$ (2) The curvature tensor on $``$ is introduced in terms of the tangent vectors $`𝑩,𝑪,𝑫`$: $$[𝑹(𝑩,𝑪)𝑫]^aR_{bcd}^aB^bC^cD^d\left(\mathrm{\Gamma }_{bd,c}^a\mathrm{\Gamma }_{bc,d}^a+\mathrm{\Gamma }_{bd}^e\mathrm{\Gamma }_{ce}^a\mathrm{\Gamma }_{bc}^e\mathrm{\Gamma }_{de}^a\right)B^bC^cD^d.$$ (3) For any two vectors $`𝑨`$ and $`𝑩`$, we define the inner product and the norm as $`𝑨𝑩=𝑨^{}𝑩𝑨^T𝑮𝑩=A^aG_{ab}B^b,`$ $`\|𝑨\|\sqrt{(𝑨𝑨)},`$ (4) respectively. Here $`𝑨^{}`$ is the cotangent vector such that $`(𝑨^{})_aA^bG_{ba}`$. We also introduce the covariant derivative $`_a`$ on $``$ acting upon a vector $`𝑨`$ as $$_bA^aA_{,b}^a+\mathrm{\Gamma }_{bc}^aA^c,$$ (5) while, the covariant derivative on $`𝑨`$ with respect to the spacetime $`x^\mu `$ is $$𝒟_\mu A^a_\mu +\mathrm{\Gamma }_{bc}^a_\mu \varphi ^bA^c.$$ (6) It should be noted that the covariant derivative reduces to the ordinary derivative when it acts upon a scalar. By varying the action (1) with respect to $`g_{\mu \nu }`$ and $`\mathit{\varphi }`$, we obtain the gravitational field equation, $$\frac{1}{\kappa _0^2}G_\nu ^\mu =T_\nu ^\mu =^\mu \mathit{\varphi }_\nu \mathit{\varphi }\delta _\nu ^\mu \left(\frac{1}{2}^\lambda \mathit{\varphi }_\lambda \mathit{\varphi }+V\right),$$ (7) and the equation of motion for the scalar fields, $$g^{\mu \nu }\left(𝒟_\mu \delta _\nu ^\lambda \mathrm{\Gamma }_{\mu \nu }^\lambda \right)_\lambda \mathit{\varphi }𝑮^1\mathbf{}^TV=0,$$ (8) where $`G_\nu ^\mu `$ and $`T_\nu ^\mu `$ are Einstein and energy-momentum tensors. It is often convenient in our analysis to represent the scalar fields as effective fluid quantities. We conclude this Section by covariantly decomposing the energy-momentum tensor into fluid quantities using a time-like four-vector $`u^\mu `$ normalized as $`u^\mu u_\mu =1`$: $`T_{\alpha \beta }=\mu u_\alpha u_\beta +ph_{\alpha \beta }+q_\alpha u_\beta +q_\beta u_\alpha +\pi _{\alpha \beta },`$ $`\mu T_{\alpha \beta }u^\alpha u^\beta ,p{\displaystyle \frac{1}{3}}T_{\alpha \beta }h^{\alpha \beta },q_\alpha T_{\beta \gamma }u^\beta h_\alpha ^\gamma ,`$ $`\pi _{\alpha \beta }T_{\gamma \delta }h_\alpha ^\gamma h_\beta ^\delta ph_{\alpha \beta }.`$ (9) Here $`\mu `$, $`p`$, $`q_\alpha `$, and $`\pi _{\alpha \beta }`$ are the energy density, pressure, energy flux, and anisotropic pressure, respectively; $`h_{\alpha \beta }g_{\alpha \beta }+u_\alpha u_\beta `$ is a projection tensor based on $`u_\alpha `$ vector, $`q_\alpha u^\alpha =0=\pi _{\alpha \beta }`$, and $`\pi _\alpha ^\alpha =0`$. The decomposition given above is in the most general form. Indeed, for a multicomponent scalar field, we have $`\mu ={\displaystyle \frac{1}{2}}\|\dot{\mathit{\varphi }}\|^2+V,p={\displaystyle \frac{1}{2}}\|\dot{\mathit{\varphi }}\|^2V,`$ $`q_\alpha =0=\pi _{\alpha \beta }.`$ (10) Equations (7) and (8) with (10) provide the fundamental expressions required for describing cosmological inflation. ### II.2 Basis vectors We continue our discussion of fields on a manifold by introducing a set of basis vectors Tent , which will prove to be useful in our analysis. First note that an arbitrary tangent vector $`𝑨`$ on $``$ can always be expanded in terms of a set of basis vectors $`\{𝒆_{(a)}\}`$ as $`𝑨A^{(a)}𝒆_{(a)}`$ with $`𝒆_{(a)}𝒆_{(b)}=G_{ab}`$. However, a different set of basis vectors generated using Gram-Schmidt orthonormalization turns out to be more convenient. From the vector $`\mathit{\varphi }`$ we can construct a set of $`N`$ linearly independent vectors $`\{\mathit{\varphi }^{(1)},\mathit{\varphi }^{(2)},\mathrm{},\mathit{\varphi }^{(N)}\}`$, where, $`\mathit{\varphi }^{(1)}\dot{\mathit{\varphi }},\mathit{\varphi }^{(n)}𝒟_t^{(n1)}\dot{\mathit{\varphi }}(n2).`$ (11) Let $`𝒆_1=\mathit{\varphi }^{(1)}/\|\mathit{\varphi }^{(1)}\|`$ be the first unit vector along the direction of the field velocity $`\dot{\mathit{\varphi }}`$. Define the second unit vector $`𝒆_2`$ to be along that part of the direction of the field accelaration $`𝒟_t\dot{\mathit{\varphi }}`$ which is normal to $`𝒆_1`$: $$𝒆_2=\frac{\mathit{\varphi }^{(2)}(𝒆_1\mathit{\varphi }^{(2)})𝒆_1}{\|\mathit{\varphi }^{(2)}(𝒆_1\mathit{\varphi }^{(2)})𝒆_1\|}.$$ (12) It is obvious from Eq. (12) that $`𝒆_1𝒆_2=0`$ by construction. A repetitive application of the Gram-Schmidt procedure then generates a set of mutually orthonormal vectors $`\{𝒆_n\}`$, which span the same subspace as the vectors $`\{\mathit{\varphi }^{(n)}\}`$. Introducing the projection operators $`𝑷_n`$ and $`𝑷_n^{}`$, which project on $`𝒆_n`$ and on the subspace perpendicular to $`\{𝒆_1,\mathrm{},𝒆_n\}`$ respectively, we may then write a general unit vector $`𝒆_n`$ as, $$𝒆_n=\frac{𝑷_{n1}^{}\mathit{\varphi }^{(n)}}{\|𝑷_{n1}^{}\mathit{\varphi }^{(n)}\|},$$ (13) where, $`𝑷_n=𝒆_n𝒆_n^{},𝑷_n^{}=𝟙{\displaystyle \underset{𝕢=\mathrm{𝟙}}{\overset{𝕟}{}}}_𝕢,_\mathrm{𝟘}^{}𝟙,`$ (14) and we define, $`𝑷^{}𝑷_1=𝒆_1𝒆_1^{},𝑷^{}𝑷_1^{}=𝟙^{}.`$ (15) Note that when the denominator in Eq. (13) vanishes, the corresponding basis vector does not exist. Using the fact that $`𝑷^{}+𝑷^{}𝟙`$, we can decompose any vector $`𝑨`$ in directions parallel and perpendicular to the field velocity: $`𝑨`$ $`=`$ $`𝑨^{}+𝑨^{}(𝑷^{}+𝑷^{})𝑨`$ (16) $`=`$ $`𝒆_1(𝒆_1𝑨)+𝒆_2(𝒆_2𝑨).`$ For the special case of just one field, $`𝒆_1`$ by definition simply reduces to the normalized scalar $`\varphi ^{(1)}/\|\varphi ^{(1)}\|`$. Hence, from Eq. (12), $`𝒆_2`$ vanishes identically, and so do all other basis vectors. Thus the decomposition (16) enables us to distinguish between single-field contributions, where only $`𝒆_1`$ survives, from multiple-field ones. ## III The perturbed universe ### III.1 Metric perturbations We know that the observed Universe is not perfectly homogeneous and isotropic. Assuming that the inhomogeneities are small enough, we can then treat the deviations using perturbation theory. In this paper we consider linear perturbations of the homogeneous and isotropic cosmological space-time described by the Friedmann-Robertson-Walker ( FRW ) model. We choose the line-element to be $`ds^2`$ $`=`$ $`a^2\left(1+2A\right)d\eta ^22a^2B_idx^i`$ (17) $`+a^2\left(g_{ij}^{(3)}+2C_{ij}\right)dx^idx^j,`$ where $`a(t)`$ is the scale factor, $`dtad\eta `$, and indices $`i,j,\mathrm{}`$, run from $`1`$ to $`3`$ labelling the spatial components. The perturbed order variables $`A(t,𝒙)`$, $`B_i(t,𝒙)`$, and $`C_{ij}(t,𝒙)`$ are based on the metric $`g_{ij}^{(3)}`$ of the 3-surfaces of constant curvature $`K=0,\pm 1`$. These are general functions of space-time which characterize the linear cosmological perturbations. Here $`t`$ and $`\eta `$ are the comoving and conformal times respectively. We denote a derivative with respect to comoving time by $`\dot{}_t`$ and one with respect to conformal time by$`_\eta `$. The Hubble parameters in terms of comoving and conformal times are defined as $`H=\dot{a}/a`$ and $`=a^{}/a=aH`$. As is evident from Eq. (17), the metric is decomposed into a background part, plus a perturbation. Correspondingly we can decompose the scalar field as $$\mathit{\varphi }(t,𝒙)=\overline{\mathit{\varphi }}(t)+𝜹\mathit{\varphi }(t,𝒙),$$ (18) where the perturbation $`𝜹\mathit{\varphi }(\delta \varphi ^a)`$ is a tangent vector on $``$, while the energy-momentum tensor is decomposed as $`T_{\mathrm{\hspace{0.33em}0}}^0=\mu (\overline{\mu }+\delta \mu ),`$ $`T_i^0={\displaystyle \frac{1}{a}}[q_i+(\mu +p)u_i](\mu +p)v_i,`$ $`T_j^i=p\delta _j^i+\pi _j^i(\overline{p}+\delta p)\delta _j^i+\pi _j^{(3)i}.`$ (19) The barred entities denote background variables. For notational simplicity we shall ignore the overbars unless required. In Eq. (19), $`v_i`$ is the frame-independent flux variable, and $`v_i`$, $`\pi _j^{(3)i}`$ are based on $`g_{ij}^{(3)}`$. From Eqs. (7) and (8), the equations for the background can be written as $`H^2={\displaystyle \frac{1}{3}}\kappa _0^2\mu {\displaystyle \frac{K}{a^2}}={\displaystyle \frac{1}{3}}\kappa _0^2\left({\displaystyle \frac{1}{2}}\|\dot{\mathit{\varphi }}\|^2+V\right){\displaystyle \frac{K}{a^2}},`$ (20) $`\dot{H}={\displaystyle \frac{1}{2}}\kappa _0^2\left(\mu +p\right)+{\displaystyle \frac{K}{a^2}}={\displaystyle \frac{1}{2}}\kappa _0^2\|\dot{\mathit{\varphi }}\|^2+{\displaystyle \frac{K}{a^2}},`$ (21) $`R=6\left(2H^2+\dot{H}+{\displaystyle \frac{K}{a^2}}\right),`$ (22) $`𝒟_t\dot{\mathit{\varphi }}+3H\dot{\mathit{\varphi }}+𝑮^1\mathbf{}^TV=0,`$ (23) $`\dot{\mu }+3H\left(\mu +p\right)=0.`$ (24) Taking the $`G_0^0`$ and $`G_i^i3G_0^0`$ components of Eq. (7) yield Eqs. (20) and (21) respectively, while Eq. (8) leads to Eq. (23). Equation (24) follows from the conservation of the energy-momentum tensor. We shall ignore the cosmological constant $`\mathrm{\Lambda }`$ in our work; nevertheless it can be easily included by making the replacements $`\mu \mu +\mathrm{\Lambda }/\kappa _0^2`$ and $`pp\mathrm{\Lambda }/\kappa _0^2`$. Note that we have explicitly retained $`K(=0,\pm 1)`$, and only at a later stage shall we set $`K=0`$. ### III.2 Scalar, vector and tensor decompositions In order to make further progress, it is customary to decompose the perturbed order variables into scalar-, vector-, and tensor-type perturbations. To the linear order, each of these three perturbations decouple from one another and evolve independently. Accordingly, the metric perurbation variables $`A(t,𝒙)`$, $`B_i(t,𝒙)`$, and $`C_{ij}(t,𝒙)`$ may be decomposed as $`A\alpha ,B_i\beta _i+B_i^{(v)},`$ $`C_{ij}g_{ij}^{(3)}\phi +\gamma _{,i\|j}+C_{(i\|j)}^{(v)}+C_{ij}^{(t)}.`$ (25) In this and the following, the superscripts $`(s)`$, $`(v)`$ and $`(t)`$ will indicate the scalar-, vector- and tensor-type perturbed order variables. The vertical bar represents a covariant derivative with respect to $`g_{ij}^{(3)}`$ and the round brackets in the subscript imply symmetrization of the indices. The scalar metric perturbations are then given by $`\alpha `$, $`\beta `$, $`\gamma `$ and $`\phi `$. The transverse-type vector perturbations $`B_i^{(v)}`$ and $`C_i^{(v)}`$ satisfy $`B_{\|i}^{(v)i}=0=C_{\|i}^{(v)i}`$ while the tensor-type perturbation $`C_{ij}^{(t)}`$ is transverse-traceless $`(C_i^{(t)i}=0=C_{i\|j}^{(t)j})`$. Both the vector and tensor perturbed order variables are based on $`g_{ij}^{(3)}`$. We define $`\mathrm{\Delta }`$ as a comoving three-space Laplacian, and introduce the following combinations of the metric variables, $`\chi a(\beta +a\dot{\gamma }),\kappa 3(H\alpha \dot{\phi }{\displaystyle \frac{\mathrm{\Delta }}{a^2}}\chi ),`$ $`\mathrm{\Psi }^{(v)}B^{(v)}+a\dot{C}^{(v)}.`$ (26) It is convenient to separate the temporal and spatial aspects of the perturbed order variables by expanding them in terms of harmonic eigenfunctions $`𝒬^{(s,v,t)}(𝒌;𝒙)`$ of the generalized Helmholtz equation Bardeen80 ; Kodamasasaki : $`𝒬_{,i}^{(s)\|i}k^2𝒬^{(s)},𝒬_i^{(s)}{\displaystyle \frac{1}{k}}𝒬_{,i}^{(s)},`$ $`𝒬^{(s)}{\displaystyle \frac{1}{k^2}}𝒬_{,i\|j}^{(s)}+{\displaystyle \frac{1}{3}}g_{ij}^{(3)}𝒬^{(s)},`$ $`𝒬_{i,j}^{(v)\|i}k^2𝒬_i^{(v)},𝒬_{ij}^{(v)}{\displaystyle \frac{1}{k}}𝒬_{(i\|j)}^{(v)},𝒬_i^{(v)\|i}0,`$ $`𝒬_{ij,k}^{(t)\|k}k^2𝒬_{ij}^{(t)},𝒬_{ij}^{(t)}𝒬_{ji}^{(t)},`$ $`𝒬_{ij}^{(t)\|j}0𝒬_i^{(t)i}.`$ (27) Here $`𝒌`$ is the wave vector in Fourier space and $`k=\|𝒌\|`$. We can then write the scalar-type perturbed order variables as $`\alpha (t,𝒙)\alpha (t,𝒌)𝒬^{(s)}(𝒌;𝒙)`$, with similar expressions for $`\beta `$, $`\gamma `$ and $`\phi `$. The vector- and tensor-type perturbations are expanded as $`B_i^{(v)}B^{(v)}𝒬_i^{(v)}`$, $`C_i^{(v)}C^{(v)}𝒬_i^{(v)}`$, and $`C_{ij}^{(t)}C^{(t)}𝒬_{ij}^{(t)}`$. In each of these harmonic expansions, a summation over the modes of the eigenfunctions is implied. In particular, the perturbed scalar fields have the expansion $$𝜹\mathit{\varphi }(t,𝒙)𝜹\mathit{\varphi }(t,𝒌)𝒬^{(s)}(𝒌;𝒙).$$ (28) From Eq. (28) we derive the important conclusion that the scalar fields on their own can generate only scalar-type perturbations. We also note that to the linear order in perturbations, the form of the equations in configuration space are identical to the corresponding ones in Fourier space. Consequently, for maintaining notational ease, we do not distinguish between the two. In a similar spirit, the fluid variables $`v_i`$ and $`\pi _j^{(3)i}`$ can be expanded in terms of the harmonics as $`v_iv^{(s)}𝒬_i^{(s)}+v^{(v)}𝒬_i^{(v)},`$ $`\pi _j^{(3)i}\pi ^{(s)}𝒬_j^{(s)i}+\pi ^{(v)}𝒬_j^{(v)i}+\pi ^{(t)}𝒬_j^{(t)i};`$ (29) while the energy-momentum tensor in Eq. (19) has the expansion $`T_{\mathrm{\hspace{0.33em}0}}^0=\mu (\overline{\mu }+\delta \mu ),`$ $`T_i^0={\displaystyle \frac{1}{k}}(\mu +p)v_{,i}^{(s)}+(\mu +p)v^{(v)}𝒬_i^{(v)},`$ $`T_j^i=(\overline{p}+\delta p)\delta _j^i+\pi ^{(s)}𝒬_j^{(s)i}+\pi ^{(v)}𝒬_j^{(v)i}+\pi ^{(t)}𝒬_j^{(t)i}.`$ (30) For a Universe having the matter sector composed exclusively of scalar fields, the quantity $`\pi _j^{(3)i}`$ in Eq. (19) vanishes identically. We then have to the perturbed order, $`\delta \mu =\dot{\mathit{\varphi }}𝒟_t\delta \mathit{\varphi }\alpha \|\dot{\mathit{\varphi }}\|^2+\mathbf{}V\delta \mathit{\varphi },`$ (31) $`\delta p=\dot{\mathit{\varphi }}𝒟_t\delta \mathit{\varphi }\alpha \|\dot{\mathit{\varphi }}\|^2\mathbf{}V\delta \mathit{\varphi },`$ (32) $`(\mu +p)v{\displaystyle \frac{a}{k}}=\dot{\mathit{\varphi }}\delta \mathit{\varphi },`$ (33) where we have written $`vv^{(s)}`$ for simplcity. It is also convenient to decompose $`\delta p`$ into an adiabatic part $`c_s^2\delta \mu `$, and an entropic perturbation $`e`$: $$\delta p=c_s^2\delta \mu +e,$$ (34) where $`c_s^2\dot{p}/\dot{\mu }`$ may be interpreted as an effective sound velocity. ### III.3 Gauge transformations As mentioned previously, our Universe shows departures from ideal homogeneity and isotropy. When the deviations are small enough, one can choose a fictitious background geometry and consider perturbations about it using infinitesimal coordinate transformations. The change in correspondence between the background and perturbed space-times represented by coordinate shifts is called a gauge transformation. Now, since relativistic gravity is invariant under coordinate transformations, the perturbations are not unique. There are different ways of mapping between the background and perturbed parts. This leads to what is called gauge degrees of freedom in the context of cosmological perturbations. Below we briefly summarize the transformation properties of various quantities to the linear order in perturbations Mukhanovetal ; Bardeen80 ; Kodamasasaki . Under a coordinate shift $`\stackrel{~}{x}^\mu =x^\mu +\xi ^\mu `$, the scalars, vectors and tensors transform as $`\stackrel{~}{𝒜}(\stackrel{~}{x}^\lambda )=𝒜(x^\lambda ),\stackrel{~}{𝒜}_\mu (\stackrel{~}{x}^\lambda )={\displaystyle \frac{x^\sigma }{\stackrel{~}{x}^\mu }}𝒜_\sigma (x^\lambda ),`$ $`\stackrel{~}{𝒜}_{\mu \nu }(\stackrel{~}{x}^\lambda )={\displaystyle \frac{x^\sigma }{\stackrel{~}{x}^\mu }}{\displaystyle \frac{x^\tau }{\stackrel{~}{x}^\nu }}𝒜_{\sigma \tau }(x^\lambda ),`$ (35) so that $`\stackrel{~}{𝒜}(x^\lambda )=𝒜(x^\lambda )𝒜_{,\sigma }\xi ^\sigma ,`$ $`\stackrel{~}{𝒜}_\mu (x^\lambda )=𝒜_\mu (x^\lambda )𝒜_{\mu ,\sigma }\xi ^\sigma 𝒜_\sigma \xi _{,\mu }^\sigma ,`$ $`\stackrel{~}{𝒜}_{\mu \nu }(x^\lambda )=𝒜_{\mu \nu }(x^\lambda )𝒜_{\mu \nu ,\sigma }\xi ^\sigma 2𝒜_{\sigma (\nu }\xi _{,\mu )}^\sigma .`$ (36) Writing the temporal part of $`\xi _\mu `$ as $`\xi ^0=a^1\xi ^t`$ and decomposing the spatial part as $`\xi _i=a^1\xi _{,i}+\xi _i^{(v)}`$, where $`\xi _i^{(v)}`$ is based on $`g_{ij}^{(3)}`$ satisfying $`\xi _{\|i}^{(v)i}=0`$, we find from Eq. (36) that the metric and matter variables transform to linear order as: $`\stackrel{~}{\alpha }=\alpha \dot{\xi ^t},\stackrel{~}{\beta }=\beta {\displaystyle \frac{1}{a}}\xi ^t+a\left({\displaystyle \frac{\xi }{a}}\right)^.,`$ $`\stackrel{~}{\gamma }=\gamma {\displaystyle \frac{1}{a}}\xi ,\stackrel{~}{\phi }=\phi H\xi ^t,\stackrel{~}{\chi }=\chi \xi ^t,`$ $`\stackrel{~}{\kappa }=\kappa +\left(3\dot{H}+{\displaystyle \frac{\mathrm{\Delta }}{a^2}}\right),\stackrel{~}{v}=v{\displaystyle \frac{1}{a}}\xi ^t,`$ $`\delta \stackrel{~}{\mu }=\delta \mu \dot{\mu }\xi ^t,\delta \stackrel{~}{p}=\delta p\dot{p}\xi ^t,𝜹\stackrel{~}{\mathit{\varphi }}=𝜹\mathit{\varphi }\dot{\mathit{\varphi }}\xi ^t,`$ $`\stackrel{~}{B}_i^{(v)}=B_i^{(v)}+a\dot{\xi }_i^{(v)},\stackrel{~}{C}_i^{(v)}=C_i^{(v)}\xi _i^{(v)},`$ $`\stackrel{~}{v}^{(v)}=v^{(v)},\stackrel{~}{\mathrm{\Psi }}^{(v)}=\mathrm{\Psi }^{(v)},`$ $`\stackrel{~}{\pi }^{(s,v,t)}=\pi ^{(s,v,t)},\stackrel{~}{C}_{ij}^{(t)}=C_{ij}^{(t)}.`$ (37) It is immediately obvious from Eq. (37) that the tensor-type perturbations are gauge-invariant. For the special case of scalar-type perturbations to the linear order, fixing the temporal part $`\xi ^t`$ of the gauge transformation leads to different gauge conditions. Table 1 summarizes some common temporal gauges. From the first equation in (37), we see that the gauge transformation of $`\alpha `$ involves the term $`\dot{\xi ^t}`$. Therefore the synchronous gauge condition $`\alpha 0`$ fixes $`\xi ^t`$ upto a constant of integration, leaving a spatially varying residual gauge mode $`\xi ^t(𝒙)`$. For the remaining gauges in Table 1, the temporal gauge mode is completely determined. We conclude by giving a few examples of gauge-invariant combination of variables: $`\phi _\chi \phi H\chi ,\alpha _\chi \alpha \dot{\chi },v_\chi v{\displaystyle \frac{k}{a}}\chi ,`$ $`\delta \mu _\chi \delta \mu \dot{\mu }\chi ,\delta p_\chi \delta p\dot{p}\chi ,`$ $`𝜹\mathit{\varphi }_\chi 𝜹\mathit{\varphi }\dot{\mathit{\varphi }}\chi ,𝜹\mathit{\varphi }_\phi 𝜹\mathit{\varphi }{\displaystyle \frac{\dot{\mathit{\varphi }}}{H}}\phi {\displaystyle \frac{\dot{\mathit{\varphi }}}{H}}\phi _{𝜹\mathit{\varphi }},`$ $`\phi _v\phi {\displaystyle \frac{aH}{k}}v,\delta \mu _v\delta \mu {\displaystyle \frac{a}{k}}\dot{\mu }v.`$ (38) Thus, in the zero-shear gauge, also known as the longitudinal, or conformal Newtonian gauge, $`\chi 0`$ is the gauge condition. We then have from Eq. (38), $`\phi _\chi \phi `$, $`\alpha _\chi \alpha `$, and $`𝜹\mathit{\varphi }_\chi 𝜹\mathit{\varphi }`$. Similarly, in the uniform-curvature gauge, it follows that $`𝜹\mathit{\varphi }_\phi 𝜹\mathit{\varphi }`$ which in turn is equivalent to $`(\dot{\mathit{\varphi }}/H)\phi _{𝜹\mathit{\varphi }}`$ in the uniform-field gauge. In the notation of Mukhanovetal , our $`\alpha _\chi `$ and $`\phi _\chi `$ correspond to their $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ respectively. ## IV Scalar perturbations in multiple-field inflation ### IV.1 Perturbation equations in the gauge-ready form We shall now briefly discuss the gauge-ready approach introduced in Hwang1 ; Hwang2 ; Hwang3 ; Hwang4 . As is well known in the theory of cosmological perturbations, a judicious choice of gauge conditions often simplifies the mathematical structure of a particular problem. For example, density perturbations with hydrodynamical fluids are most conveniently treated using the comoving gauge, while the uniform-curvature gauge simplifies the analysis of perturbations due to minimally coupled scalar fields. Since, in general, we do not know the optimal gauge condition beforehand, it becomes advantageous to express the perturbations without imposing a specific temporal gauge condition. In other words, we write the governing equations in the gauge-ready form, which would give us the freedom to choose different gauge conditions, as adapted to the problem, at a later stage in the calculations. Once the temporal gauge mode is completely fixed so that no further gauge degrees of freedom are left, the resulting variables would then be gauge-invariant. Moreover, when a solution in a particular gauge is known, we can then easily derive the corresponding solution in other gauges, as well as in gauge-invariant forms. This is the basic concept of the gauge-ready method. To implement this gauge-ready strategy, it is most convenient to derive the perturbed set of equations from the (3+1) ADM ADM , and the (1+3) covariant covariant formulations of Einstein gravity. A complete set of these equations may be found in the Appendix of Ref. Hwang1 . In this Section we write the equations for scalar-type perturbations in the gauge-ready form. Definition of $`\kappa `$: $$\dot{\phi }=H\alpha \frac{1}{3}\kappa +\frac{1}{3}\frac{k^2}{a^2}\chi .$$ (39) ADM energy constraint ($`G_0^0`$ component of the field equation): $$\frac{k^23K}{a^2}\phi +H\kappa =\frac{1}{2}\kappa _0^2\delta \mu .$$ (40) ADM momentum constraint ($`G_i^0`$ component): $$\kappa \frac{k^23K}{a^2}\chi =\frac{3}{2}\kappa _0^2(\mu +p)\frac{a}{k}v.$$ (41) ADM propagation($`G_j^i\frac{1}{3}\delta _j^iG_k^k`$ component): $$\dot{\chi }+H\chi \alpha \phi =\kappa _0^2\frac{a^2}{k^2}\pi ^{(s)}.$$ (42) Raychaudhuri equation ($`G_i^iG_0^0`$ component): $$\dot{\kappa }+2H\kappa +\left(3\dot{H}\frac{k^2}{a^2}\right)\alpha =\frac{1}{2}\kappa _0^2(\delta \mu +3\delta p).$$ (43) Equation of motion for scalar fields: $`\left(𝒟_t^2+3H𝒟_t{\displaystyle \frac{\mathrm{\Delta }}{a^2}}+𝑴^2\right)𝜹\mathit{\varphi }`$ $`=`$ $`\left(\dot{\alpha }3\dot{\phi }{\displaystyle \frac{\mathrm{\Delta }}{a^2}}\chi \right)\dot{\mathit{\varphi }}`$ (44) $``$ $`2\alpha 𝑮^1\mathbf{}^TV.`$ Energy conservation: $$\delta \dot{\mu }+3H(\delta \mu +\delta p)=(\mu +p)\left(\kappa 3H\alpha \frac{k}{a}v\right).$$ (45) Momentum conservation: $`{\displaystyle \frac{[a^4(\mu +p)v]^\dot{}}{a^4(\mu +p)}}`$ $`={\displaystyle \frac{k}{a}}\left[\alpha +{\displaystyle \frac{1}{\mu +p}}\left(\delta p{\displaystyle \frac{2}{3}}{\displaystyle \frac{k^23K}{k^2}}\pi ^{(s)}\right)\right].`$ (46) In the above equations, $`\delta \mu `$ and $`\delta p`$ are given by Eqs. (31) and (32) respectively, while $$𝑴^2=𝑮^1\mathbf{}^T\mathbf{}V𝑹(\dot{\mathit{\varphi }},\dot{\mathit{\varphi }}).$$ (47) Note that these equations are valid for any $`K`$, and for a scalar field, $`\pi ^{(s)}=0`$. Equations (39)-(46), together with the background equations (20)-(24), and the perturbed order variables for the scalar fields (31)-(33), provide a complete set of equations for analyzing scalar-type cosmological perturbations with multicomponent scalar fields. As we have not chosen a specific gauge so far, Eqs. (39)-(46) are therefore in the gauge-ready form. This allows us to impose any one of the available temporal gauge conditions, which would then fix the temporal gauge mode completely, leading to gauge-invariant variables. ### IV.2 Perturbation equations using gauge-invariant variables In order to illustrate the gauge-ready method, we derive some useful expressions in terms of the gauge-invariant variables introduced in Section III.3. From Eqs. (40) and (41) we obtain $$\frac{k^23K}{a^2}\phi _\chi =\frac{1}{2}\kappa _0^2\delta \mu _v.$$ (48) Eq. (42) can be written as $$\alpha _\chi +\phi _\chi =\kappa _0^2\frac{a^2}{k^2}\pi ^{(s)}.$$ (49) Eqs. (41),(42) and (39) lead to $$\dot{\phi }_\chi H\alpha _\chi =\frac{1}{2}\kappa _0^2(\mu +p)\frac{a}{k}v_\chi ,$$ (50) Eqs. (45),(46) with (41) yield $`\delta \dot{\mu }_v+3H\delta \mu _v`$ $`={\displaystyle \frac{k^23K}{a^2}}\left[(\mu +p){\displaystyle \frac{a}{k}}v_\chi +2H{\displaystyle \frac{a^2}{k^2}}\pi ^{(s)}\right],`$ (51) while Eqs. (42) and (46) give $`\dot{v}_\chi +Hv_\chi `$ $`={\displaystyle \frac{k}{a}}\left[\alpha _\chi +{\displaystyle \frac{\delta p_v}{\mu +p}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{k^23K}{a^2}}{\displaystyle \frac{\pi ^{(s)}}{\mu +p}}\right].`$ (52) Combining Eqs. (48)-(52) we can derive $`\ddot{\phi }_\chi +(4+3c_s^2)H\phi _\chi c_s^2{\displaystyle \frac{\mathrm{\Delta }}{a^2}}\phi _\chi `$ $`+\left[(\mu c_s^2p)2(1+3c_s^2){\displaystyle \frac{K}{a^2}}\right]\phi _\chi `$ $`={\displaystyle \frac{1}{2}}\kappa _0^2\left(e{\displaystyle \frac{2}{3}}\pi ^{(s)}\right)`$ $`{\displaystyle \frac{1}{2}}\kappa _0^2{\displaystyle \frac{\mu +p}{H}}\left({\displaystyle \frac{2H^2}{\mu +p}}{\displaystyle \frac{a^2}{k^2}}\pi ^{(s)}\right),`$ (53) where we used Eq. (34). For the explicit forms of the gauge-invariant variables used in these equations, see Eq. (38). From Eq. (49) we can draw the important conclusion that, for scalar-fields, $`\alpha _\chi =\phi _\chi `$, since $`\pi ^{(s)}=0`$. Using this result the equation of motion for scalar fields becomes $`\left(𝒟_\eta ^2+2H𝒟_\eta \mathrm{\Delta }+a^2𝑴^2\right)𝜹\mathit{\varphi }_\chi =4\phi _\chi ^{}\mathit{\varphi }^{}`$ $`+2a^2\phi _\chi 𝑮^1\mathbf{}^TV.`$ (54) It is also convenient to re-write Eqs. (50) and (53) for the case of scalar fields as $`\phi _\chi ^{}+\phi _\chi ={\displaystyle \frac{1}{2}}\kappa _0^2\mathit{\varphi }^{}𝜹\mathit{\varphi }_\chi ,`$ (55) $`\phi _\chi ^{\prime \prime }+6\phi _\chi ^{}\mathrm{\Delta }\phi _\chi +2\left[^{}+2(^2K)\right]\phi _\chi `$ $`=\kappa _0^2a^2\mathbf{}V𝜹\mathit{\varphi }_\chi ,`$ (56) where we used Eq. (33), and the relations $`e=\delta pc_s^2\delta \mu =\delta p_\chi c_s^2\delta \mu _\chi ,`$ $`(1c_s^2)\delta \mu _\chi e=\delta \mu _\chi \delta p_\chi =\mathbf{}V𝜹\mathit{\varphi }_\chi .`$ (57) Eq. (55) is often called the constraint equation. Note that we have written Eqs. (54)-(56)in conformal time. These equations contain most of the physics related to inflationary cosmological perturbations. They are expressed in terms of gauge-invariant forms of the variables, and from the discussion at the end of Section III.3, we see that they retain the same algebraic forms in the zero-shear gauge. We end this Section by expressing Eq. (56) in a different way. Observe that according to Eq. (16), $`𝜹\mathit{\varphi }_\chi `$ may be decomposed into components parallel and perpendicular to the field velocity, $`𝜹\mathit{\varphi }_\chi =𝜹\mathit{\varphi }_\chi ^{}+𝜹\mathit{\varphi }_\chi ^{}`$. Using the background equation (23), the constraint equation (55), and the fact that $`\|\mathit{\varphi }^{}\|^{}\|\mathit{\varphi }^{}\|=(𝒟_\eta \mathit{\varphi }^{})\mathit{\varphi }^{}`$, we can write Eq. (56) as $`\phi _\chi ^{\prime \prime }`$ $`+`$ $`2\left({\displaystyle \frac{\|\mathit{\varphi }^{}\|^{}}{\|\mathit{\varphi }^{}\|}}\right)\phi _\chi ^{}`$ (58) $`+`$ $`2\left[\left(^{}{\displaystyle \frac{\|\mathit{\varphi }^{}\|^{}}{\|\mathit{\varphi }^{}\|}}\right)2K\right]\phi _\chi \mathrm{\Delta }\phi _\chi `$ $`=`$ $`\kappa _0^2(𝒟_\eta \mathit{\varphi }^{})𝜹\mathit{\varphi }^{}.`$ Following our discussion in Section II.2, we know that the perpendicular component of field perturbation vanishes when there is only one field. In this case, the right hand side of Eq. (58) vanishes, and the resulting equation is well known in the theory of single field inflationary perturbations Mukhanovetal . ## V Solutions of the perturbation equations ### V.1 Slow-roll variables To proceed further with our analysis, we make the assumption that the Universe has undergone inflation to complete flatness, so that henceforth we can set $`K=0`$. We introduce the functions, $`ϵ(\mathit{\varphi }){\displaystyle \frac{\dot{H}}{H^2}},𝜼(\mathit{\varphi }){\displaystyle \frac{\mathit{\varphi }^{(2)}}{H\|\dot{\mathit{\varphi }}\|}},`$ (59) known as the slow-roll variables. Using Eq. (16), $`𝜼`$ is decomposed into parallel and perpendicular components: $$\eta ^{}=𝒆_1𝜼=\frac{𝒟_t\dot{\mathit{\varphi }}\dot{\mathit{\varphi }}}{H\|\dot{\mathit{\varphi }}\|^2},\eta ^{}=𝒆_2𝜼=\frac{\|(𝒟_t\dot{\mathit{\varphi }})^{}\|}{H\|\dot{\mathit{\varphi }}\|}.$$ (60) The standard slow-roll assumptions are $$ϵ=O(\zeta ),\eta ^{}=O(\zeta ),\eta ^{}=O(\zeta ),$$ (61) for some small parameter $`\zeta `$, with $`ϵ`$, $`\sqrt{ϵ}\eta ^{}`$ and $`\sqrt{ϵ}\eta ^{}`$ much smaller than unity. If in an expansion in slow-roll variables we neglect terms of order $`O(\zeta ^2)`$, we claim that expansion to be of first order in slow-roll. Thus terms with $`ϵ^2`$, $`ϵ\eta ^{}`$, etc. are of second order. Note that the definitions (59) remain valid irrespective of the slow-roll assumptions. We present some useful relations involving the slow-roll variables: $`^{}=^2(1ϵ),{\displaystyle \frac{\|\mathit{\varphi }^{}\|^{}}{\|\mathit{\varphi }^{}\|}}=(1+\eta ^{}),`$ $`𝒟_\eta \mathit{\varphi }^{}=\|\mathit{\varphi }^{}\|(𝜼+𝒆_1)=\kappa _0^1\sqrt{2}^2\sqrt{ϵ}(𝜼+𝒆_1),`$ $`^2ϵ={\displaystyle \frac{1}{2}}\kappa _0^2\|\mathit{\varphi }^{}\|^2,ϵ^{}=2ϵ(ϵ+\eta ^{}).`$ (62) ### V.2 Analysis using gauge-invariant variables In order to solve the system of perturbation equations (54), (55) and (58), we shall find it convenient to introduce the variables $`𝒒=a\left(𝜹\mathit{\varphi }_\chi {\displaystyle \frac{\mathit{\varphi }^{}}{}}\phi _\chi \right)=a\left(𝜹\mathit{\varphi }{\displaystyle \frac{\mathit{\varphi }^{}}{}}\phi \right),`$ (63) $`u={\displaystyle \frac{a}{\|\mathit{\varphi }^{}\|}}\phi _\chi ,`$ (64) where the second equality for $`𝒒`$ in Eq. (63) follows from Eq. (38). Indeed, $`𝒒`$ is gauge-invariant, and is a natural generalization of the single field Sasaki-Mukhanov variable Mukhanovetal . Using the slow-roll variables (59), together with some of the relations (62), the constraint equation (55) can be written in terms of $`𝒒`$ as $$\phi _\chi ^{}+(1+ϵ)\phi _\chi =\frac{1}{2}\kappa _0^2\mathit{\varphi }^{}\frac{𝒒}{a},$$ (65) From Eq. (54), the scalar field perturbations satisfy $$𝒟_\eta ^2𝒒(\mathrm{\Delta }^2𝛀)𝒒=0,$$ (66) where $$𝛀=\frac{a^2𝑴^2}{^2}(2ϵ)𝟙\mathrm{𝟚}ϵ\left((\mathrm{𝟛}+ϵ)^{}+𝕖_\mathrm{𝟙}𝜼^{}+𝜼𝕖_\mathrm{𝟙}^{}\right),$$ (67) and use has also been made of Eq. (65). The corresponding Lagrangean $``$ follows from Eq. (66): $`S`$ $`=`$ $`{\displaystyle \sqrt{g^{(3)}}𝑑\eta d^3𝒙}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \left(𝒟_\eta 𝒒^{}𝒟_\eta 𝒒+𝒒^{}(\mathrm{\Delta }^2𝛀)𝒒\right)\sqrt{g^{(3)}}𝑑\eta d^3𝒙}.`$ Here $`g^{(3)}`$ is the determinant of the metric $`g_{ij}^{(3)}`$ of the 3-surfaces of constant curvature $`K=0`$, see below Eq. (17). The equation of motion for $`u`$ is obtained by substituting its definition (64) into Eq. (58): $`u^{\prime \prime }\mathrm{\Delta }u{\displaystyle \frac{\theta ^{\prime \prime }}{\theta }}u=\kappa _0^2\eta ^{}q_2,q_2𝒆_2𝒒,`$ $`\theta {\displaystyle \frac{}{a\|\mathit{\varphi }^{}\|}}={\displaystyle \frac{\kappa _0}{\sqrt{2}}}{\displaystyle \frac{1}{a\sqrt{ϵ}}}.`$ (69) For later use, we also express Eq. (65) in terms of $`u`$ and $`q`$ as $$u^{}+\frac{(1/\theta )^{}}{1/\theta }u=\frac{1}{2}q_1,q_1𝒆_1𝒒.$$ (70) Differentiating Eq. (70) once with respect to the conformal time and using Eq. (69), we obtain the relation $$\frac{1}{2}\left(q_1^{}\frac{(1/\theta )^{}}{1/\theta }q_1\right)\kappa _0^2\eta ^{}q_2=\mathrm{\Delta }u.$$ (71) Although the equations (65),(66) and (69) have been expressed in terms of the slow-roll variables, they are exact, and no slow-roll approximation has yet been made. Observe that, to the leading order in slow-roll, the perturbation variables $`𝒒`$ and $`u`$ decouple, whereas at first order, mixing between these occur. ### V.3 Quantization of the scalar perturbations We now study the quantization of the density perturbations described by the Lagrangean in Eq. (LABEL:fieldaction). We start by introducing the matrix $`Z_{mn}`$ defined as $$(Z)_{mn}=(Z^T)_{mn}=\frac{1}{}𝒆_m𝒟_\eta 𝒆_n,$$ (72) where the second equality follows from $`𝒟_\eta (𝒆_m𝒆_n)=0`$. Thus $`Z`$ is antisymmetric and traceless $`(\text{Tr}Z=0)`$. Upon expanding $`𝒒=q_m𝒆_m`$, using the basis $`\{𝒆_m\}`$, it follows from Eq. (LABEL:fieldaction), $$=\frac{1}{2}(q^{}+Zq)^T(q^{}+Zq)+\frac{1}{2}q^T(\mathrm{\Delta }^2\mathrm{\Omega })q,$$ (73) where $`(\mathrm{\Omega })_{mn}=𝒆_m^{}𝛀𝒆_n`$, and for notational ease, we have suppresed the indices $`m,n`$. It will prove convenient to reduce the Lagrangean (73) to the canonical form. We redefine $`q`$ using a new matrix $`R`$ as $$q(\eta )=R(\eta )Q(\eta ),R^{}+ZR=0,\stackrel{~}{\mathrm{\Omega }}=R^T\mathrm{\Omega }R.$$ (74) From the equation of motion (74) for $`R`$, it follows that $`R^TR`$ and $`\text{det}R`$ are constants, so that $`R`$ represents a rotation. Without any loss of generality, the initial value of $`R`$ may be chosen as $`R(\eta _0)=𝟙`$. Substituting the variables defined in Eq. (74) into Eq. (73) yields $$=\frac{1}{2}Q^TQ^{}+\frac{1}{2}Q^T(\mathrm{\Delta }^2\stackrel{~}{\mathrm{\Omega }})Q.$$ (75) To proceed with the quantization, we employ the canonical quantization procedure to the Lagrangean (75). The momentum $`\mathrm{\Pi }`$ canonically conjugate to $`Q`$ is $$\mathrm{\Pi }(\eta ,𝒙)=/Q^T=Q^{}(\eta ,𝒙).$$ (76) The Hamiltonian is then given by $$=\frac{1}{2}\mathrm{\Pi }^T\mathrm{\Pi }\frac{1}{2}Q^T(\mathrm{\Delta }^2\stackrel{~}{\mathrm{\Omega }})Q.$$ (77) The canonically conjugate variables $`(Q,\mathrm{\Pi })`$ are promoted to quantum operators $`(\widehat{Q},\widehat{\mathrm{\Pi }})`$ satisfying the commutation relations $`[\alpha ^T\widehat{Q}(\eta ,𝒙),\beta \widehat{Q}(\eta ,𝒙^{\mathbf{}})]=[\alpha ^T\widehat{\mathrm{\Pi }}(\eta ,𝒙),\beta \widehat{\mathrm{\Pi }}(\eta ,𝒙^{\mathbf{}})]=0,`$ $`[\alpha ^T\widehat{Q}(\eta ,𝒙),\beta \widehat{\mathrm{\Pi }}(\eta ,𝒙^{\mathbf{}})]=i\alpha ^T\beta \delta (𝒙𝒙^{\mathbf{}}),`$ (78) where the delta function is normalized as $$\delta (𝒙𝒙^{\mathbf{}})\sqrt{g^{(3)}}d^3𝒙,$$ (79) and we have introduced the vectors $`\alpha ,\beta `$ with components $`\alpha _m,\beta _m`$ in the basis $`\{𝒆_m\}`$ to avoid writing the indices $`m,n`$ in the commutators. Since we are considering spatially flat hypersurfaces $`(K=0)`$, the operator $`\widehat{Q}`$ may be expanded in a plane wave basis as $$\widehat{Q}=\frac{d^3𝒌}{(2\pi )^{3/2}}\left[Q_k^{}(\eta )\widehat{a}_𝒌e^{i𝒌x}+Q_k(\eta )\widehat{a}_𝒌^{}e^{i𝒌x}\right],$$ (80) with a similar expansion for $`\widehat{\mathrm{\Pi }}`$. It immediately follows from Eq. (74) that $`q`$ must now be interpreted as the operator $`\widehat{q}`$ with modes $$q_k(\eta )=R(\eta )Q_k(\eta ),$$ (81) satisfying a mode expansion identical to Eq. (80). The creation and annihilitation operators $`\widehat{a}_𝒌^{}`$ and $`\widehat{a}_𝒌`$ satisfy $`[\alpha ^T\widehat{a}_𝒌,\beta \widehat{a}_𝒌^{}]=[\alpha ^T\widehat{a}_𝒌^{},\beta \widehat{a}_𝒌^{}^{}]=0,`$ $`[\alpha ^T\widehat{a}_𝒌,\beta \widehat{a}_𝒌^{}^{}]=\alpha ^T\beta \delta (𝒌𝒌^{}).`$ (82) In order that the commutation relations (78) and (82) be consistent, the following Wronskian condition must be satisfied, $$W\{Q_k,Q_k^{}\}Q_k^{}(\eta )Q_k^{}(\eta )Q_k^{}(\eta )Q_k(\eta )=i.$$ (83) From the mode expansion (80) and the Hamiltonian (77), it follows that the equation of motion for $`Q_k`$ is $$Q_k^{\prime \prime }+(k^2+^2\stackrel{~}{\mathrm{\Omega }})Q_k=0.$$ (84) It may be easily verified using Eq. (84) that the Wronskian satisfies $`dW\{Q_k,Q_k^{}\}/d\eta =0`$. We also interpret the variable $`u`$ introduced in Eq. (64) as an operator $`\widehat{u}`$, and after performing a mode expansion identical to that of $`\widehat{Q}`$ in Eq. (80), it follows from Eq. (69) that the modes $`u_k`$ satisfy $$u_k^{\prime \prime }+\left(k^2\frac{\theta ^{\prime \prime }}{\theta }\right)u_k=\kappa _0^2\eta ^{}q_{2k},q_{2k}(𝒆_2𝒆_m)q_k,$$ (85) or, equivalently, from Eq. (71), $`\kappa _0^2\eta ^{}q_{2k}{\displaystyle \frac{1}{2}}\left(q_{1k}^{}{\displaystyle \frac{(1/\theta )^{}}{1/\theta }}q_{1k}\right)=k^2u_k,`$ $`q_{1k}(𝒆_1𝒆_m)q_k.`$ (86) ### V.4 First order solution As a prelude to presenting the solution of the perturbation equations to the first order in slow-roll, we briefly discuss the issue of the ambiguity in the choice of the vacuum state when quantizing fields in an expanding FRW background birrell . In ordinary Minkowski space-time, there exists a unique time direction, as well as distinct time-invariant positive- and negative-frequency modes. However, when quantizing in a curved background, there is neither a distinct time direction, nor a notion of time-invariant mode. Consequently, there is no unique vacuum state either. Hence if we have a positive-frequency mode $`Q_k^+`$ at times $`\eta <\eta _0<0`$, with the initial vacuum state $`\widehat{a}_𝒌\|0=0`$, then at later times $`\eta >\|\eta _0\|`$, the modes will in general be described by a linear superposition of positive- and negative-frequency modes $`Q_k^+,Q_k^{}`$, related by means of a Bogoliubov transformation: $`Q_k(\eta )=\lambda _k(\eta )Q_k^+(\eta _0)+\mu _k(\eta )Q_k^{}(\eta _0),`$ $`\|\lambda _k(\eta )\|^2\|\mu _k(\eta )\|^2=1.`$ (87) For the coefficients $`\lambda _k`$ and $`\mu _k`$, one often makes the choice of the initial values as $$\lambda _k(\eta _0)=1,\mu _k(\eta _0)=0,$$ (88) for the adiabatic vacuum (or the Bunch-Davies vacuum in de-Sitter space), which corresponds to the positive-frequency solution in the Minkowski space. To proceed further, it is convenient to introduce the time $`\eta _{}`$ when the mode with wave number $`k`$ crosses the Hubble radius during inflation, so that the relation $$(\eta _{})=k$$ (89) is satisfied for each $`k`$. Consequently, the inflationary epoch can be separated into three regions: the sub-horizon region ($`k`$), the transition region ($`k`$), and the super-horizon region ($`k`$). We now discuss each of these in turn. In the sub-horizon region, we solve Eq. (84) with the $`^2\stackrel{~}{\mathrm{\Omega }}`$ term subdominant compared to $`k^2`$. The solution is obtained in the limit $`k/\mathrm{}`$ for fixed $`k`$ as $$Q_k(\eta )=\frac{1}{\sqrt{2k}}e^{ik(\eta \eta _0)},R(\eta _0)=𝟙.$$ (90) Since one is usually interested in calculating quantities at the end of inflation, this region is therefore irrelevant. We consider next the transition region. It will prove useful to introduce the time $`\eta _{}`$ when the sub-horizon epoch ends and the transition region begins. In a sufficiently small interval around $`\eta _{}`$ we can then apply slow-roll to the Eq. (84) keeeping all the terms, but taking the slow-roll functions to be constant to the first order. The initial conditions are chosen as $$Q_k(\eta _{})=\frac{1}{\sqrt{2k}}𝟙,_𝕜^{}(\eta _{})=\frac{𝕚\sqrt{𝕜}}{\sqrt{\mathrm{𝟚}}}𝟙,(\eta _{})=𝟙.$$ (91) Integrating the relation for $`^{}`$ in Eq. (62) with the initial conditions for the transition region, we obtain $$(\eta )=\frac{1}{(1ϵ_{})\eta },\eta _{}=\frac{1}{(1ϵ_{})k},$$ (92) so that $`(\eta _{})=k`$. Differentiating $`\theta `$ in Eq. (69) yields $`\theta ^{}/\theta =(1+ϵ+\eta ^{})`$, which can be integrated with the result $`\theta (z)=\theta _{}\left({\displaystyle \frac{z}{z_{}}}\right)^{(1+2ϵ_{}+\eta _{}^{})},zk\eta ,`$ $`\theta _{}={\displaystyle \frac{\kappa }{\sqrt{2}}}{\displaystyle \frac{H_{}}{k\sqrt{ϵ}_{}}},z_{}k\eta _{}`$ (93) The differential equation (74) for the rotation matrix $`R`$ can be solved with the initial conditions (91) leading to $$R(z)=\left(\frac{z}{z_{}}\right)^{(1ϵ_{})^1Z_{}},z_{}k\eta _{}.$$ (94) Since the time-dependent terms in the matrix $`\mathrm{\Omega }`$ in Eq. (67) are of first order, we can take $`\mathrm{\Omega }=\mathrm{\Omega }(\eta _{})\mathrm{\Omega }_{}`$ in the transition region. Then the matrix $`\stackrel{~}{\mathrm{\Omega }}`$ is given by $`\stackrel{~}{\mathrm{\Omega }}`$ $`=`$ $`R^1(z)\mathrm{\Omega }_{}R(z)=\mathrm{\Omega }_{}[\mathrm{\Omega }_{},Z_{}]\text{ln}{\displaystyle \frac{z}{z_{}}}`$ (95) $`=`$ $`\mathrm{\Omega }_{}+3[\delta _{},Z_{}]\left(\text{ln}{\displaystyle \frac{z}{z_{}}}+{\displaystyle \frac{3}{4}}\text{ln}ϵ_{}\right),`$ with $`\delta (\eta )`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(2𝟙+{\displaystyle \frac{\mathbb{\Omega }}{(\mathrm{𝟙}ϵ)^\mathrm{𝟚}}}\right)`$ $`=`$ $`ϵ𝟙{\displaystyle \frac{𝕒^\mathrm{𝟚}𝕄^\mathrm{𝟚}}{\mathrm{𝟛}^\mathrm{𝟚}}}+\mathrm{𝟚}ϵ(𝕖_\mathrm{𝟙}𝕖_𝕞)(𝕖_\mathrm{𝟙}𝕖_𝕟)^𝕋,`$ $`\delta _{}`$ $`=`$ $`\delta (\eta _{}).`$ (96) Here the second equality is valid to the first order in slow-roll, and we have used the notation $`M^2𝒆_m^{}𝑴^2𝒆_n`$. We also made the assumption that those components of $`a^2M^2/^2`$ which cannot be expressed in terms of the slow-roll variables are of first order. Because $`\delta _{}`$ and $`Z_{}`$ are both of first order, we can take $`\stackrel{~}{\mathrm{\Omega }}=\mathrm{\Omega }_{}`$ in Eq. (95) to be a first order quantity. In order to write the equation for the mode $`Q_k`$ in the transition region, we will find it convenient to define $`\overline{Q}_kR_{}Q_k(z)`$ and $`\overline{\mathrm{\Omega }}=R_{}\stackrel{~}{\mathrm{\Omega }}R_{}^1`$, with $`R_{}R(z_{})`$. From Eq. (94), we have to the first order, $`Q_k(z)=\overline{Q}_k(z)`$, while from Eq. (95) we conclude that $`\overline{\mathrm{\Omega }}=\mathrm{\Omega }_{}`$ within a small region around $`z_{}`$. Using the above results in Eq. (84), the mode equation for $`Q_k`$ may be written in terms of $`\overline{Q}_k`$ as $$\overline{Q}_{k,zz}+\left(𝟙\frac{\nu _{}^\mathrm{𝟚}\frac{\mathrm{𝟙}}{\mathrm{𝟜}}}{𝕫^\mathrm{𝟚}}\right)\overline{Q}_k=0,\nu _{}^2=\frac{9}{4}𝟙+\mathrm{𝟛}\delta _{}.$$ (97) This equation is similar to the one obtained for the single-field inflation, except that this is a matrix equation. The solution is then given in terms of the Hankel functions of matrix valued order $`\nu _{}`$, $`\overline{Q}_k(z)=\sqrt{z}[c_1(k)H_\nu _{}^{(1)}(z)+c_2(k)H_\nu _{}^{(2)}(z)],`$ $`\nu _{}={\displaystyle \frac{3}{2}}𝟙+\delta _{}.`$ (98) We wish to match the solution in Eq. (98) so that in the limit $`k/\mathrm{}`$, the modes approach plane waves, $`\overline{Q}_k(z)=e^{iz}/\sqrt{2k}`$, see (90). For $`\|z\|1`$, the Hankel functions have the asymptotic forms, $`H_\nu _{}^{(1)}(z)\sqrt{2/(\pi z)}e^{i\{z(\nu _{}+1/2)\pi /2\}},`$ $`H_\nu _{}^{(2)}(z)\sqrt{2/(\pi z)}e^{i\{z(\nu _{}+1/2)\pi /2\}}.`$ (99) We set $`c_1(k)=\sqrt{\pi /(4k)}e^{i(\nu _{}+1/2)\pi /2}`$, and $`c_2(k)=0`$. The phase factor of $`c_1(k)`$ is chosen in order to match with Eq. (90) at short scales, while the factor of $`\sqrt{\pi /(4k)}`$ ensures conformity with the Wronskian in Eq. (83). Therefore the final solution with the appropriate normalization is $$\overline{Q}_k(z)=\sqrt{\pi /(4k)}e^{i(\nu _{}+1/2)\pi /2}\sqrt{z}H_\nu _{}^{(1)}(z).$$ (100) It is worth mentioning that that the matrix valued Hankel functions are to be interpreted as series expansions, just like the usual Hankel functions. We finally discuss the solution in the super-horizon region. On super-horizon scales we have $`\|z\|1`$, for which the asymptotic form of the Hankel function is $$H_\nu _{}^{(1)}(z)\sqrt{2/\pi }e^{i\pi /2}2^{\nu _{}3/2}\frac{\mathrm{\Gamma }(\nu _{})}{\mathrm{\Gamma }(3/2)}z^\nu _{},$$ (101) so that the asymptotic solution for $`\overline{Q}_k(z)`$ in the super-horizon region is given by $`\overline{Q}_k(z)`$ $``$ $`(1/\sqrt{2k})e^{i(\nu _{}1/2)\pi /2}2^{\nu _{}3/2}{\displaystyle \frac{\mathrm{\Gamma }(\nu _{})}{\mathrm{\Gamma }(3/2)}}z^{\frac{1}{2}𝟙\nu _{}},`$ $``$ $`(1/\sqrt{2k})e^{i(\nu _{}1/2+2\delta _{})\pi /2}E_{}(z/z_{})^{𝟙\delta _{}},`$ where $$E_{}(1ϵ_{})𝟙+(\mathrm{𝟚}\gamma _𝔼\text{ln}\mathbb{\hspace{0.33em}2})\delta _{},$$ (103) and $`\gamma _E0.5772`$ is the Euler constant. In this region since $`k/0`$, we can also solve Eq. (85) ignoring the $`k^2`$ dependent term, leading to $`u_k(\eta )=u_{Pk}+C_k\theta +D_k\theta {\displaystyle _\eta _{}^\eta }{\displaystyle \frac{d\eta ^{}}{\theta ^2(\eta ^{})}},`$ $`u_{Pk}=\theta {\displaystyle _\eta _{}^\eta }{\displaystyle \frac{d\eta ^{}}{\theta ^2}}{\displaystyle _\eta _{}^\eta ^{}}𝑑\eta ^{\prime \prime }\theta \kappa _0^2\eta ^{}q_{2k},`$ (104) where $`C_k`$ and $`D_k`$ are constants of integration, and $`u_{Pk}`$ is a particular solution. Note that since $`\theta `$ is a rapidly decaying function, we can ignore $`C_k`$ compared to $`D_k`$. In the same approximation, the solution of Eq. (86) is $$q_{1k}=d_k(1/\theta )+2(1/\theta )_\eta _{}^\eta 𝑑\eta ^{}\theta \kappa _0^2\eta ^{}q_{2k}.$$ (105) From Eq. (70) we see that the integration constants $`D_k`$ and $`d_k`$ are related by $`D_k=\frac{1}{2}d_k`$. Considering the region where $`\eta `$ is sufficiently close to $`\eta _{}`$, the integral in Eq. (105) may then be neglected, so that using Eq. (93), we can write $`q_{1k}=2D_k(1/\theta _{})(z/z_{})^1`$. Taking into account the asymptotic solution (LABEL:qbarsuper), and the fact that $`q_k=(𝒆_1𝒆_m)^Tq_{1k}`$, we finally obtain, $$D_k=(1/2\sqrt{2k})e^{i(\nu _{}1/2+2\delta _{})\pi /2}\theta _{}(𝒆_1𝒆_m)^TE_{}.$$ (106) Thus the integration constant in Eq. (104) is completely determined to first order in slow-roll. Inserting the result (106) for $`D_k`$ in (104), and using the relation $`a_{}H_{}=k`$, we finally arrive at $$u_k=\frac{1}{(2k)^{3/2}}e^{i(\nu _{}1/2+2\delta _{})\pi /2}\frac{H_{}}{\sqrt{ϵ_{}}}\left[𝒜(t_{},t)(𝒆_1𝒆_m)^T+(t_{},t)\right]E_{},$$ (107) where we ignored $`C_k`$, and $`𝒜(t_{},t)={\displaystyle \frac{1}{a\sqrt{ϵ}}}{\displaystyle _t_{}^t}𝑑t^{}a\left({\displaystyle \frac{1}{H}}\right)^.,(t_{},t)={\displaystyle \frac{1}{a\sqrt{ϵ}}}{\displaystyle _t_{}^t}𝑑t^{}a\left({\displaystyle \frac{1}{H}}\right)^.𝒰(t_{},t),`$ $`𝒰(t_{},t)=2\kappa _0^2{\displaystyle _t_{}^t}𝑑t^{}H\eta ^{}\sqrt{{\displaystyle \frac{ϵ_{}}{ϵ}}}{\displaystyle \frac{a_{}}{a}}(𝒆_2𝒆_m)^TR{\displaystyle \frac{Q_k}{Q_k}}.`$ (108) Here $`Q_k`$ is the value of the asymptotic solution (LABEL:qbarsuper) for $`Q_k`$ evaluated at $`\eta =\eta _{}`$. Observe that the solution (107) for $`u_k`$ is expressed entirely in terms of background quantities and comoving time. This concludes our discussion of scalar perturbations in multiple field inflation. ### V.5 Vector and tensor perturbations For the sake of completeness, we now present a brief discussion of vector- and tensor-type perturbations. From the $`G_i^0`$ component of Eq. (7), together with Eq. (30), we have $$\frac{1}{2}k^2\mathrm{\Psi }^{(v)}=\kappa _0^2a^2(\mu +p)v^{(v)},$$ (109) while the condition $`T_{i;\mu }^\mu =0`$ yields $$\frac{1}{a^4}\left[a^4(\mu +p)v^{(v)}\right]^{}=\frac{1}{2}k\pi ^{(v)}.$$ (110) Equations (109) and (110) describe the vector-type, or rotational perturbations. Since vector sources are absent when the matter sector is composed entirely of scalar fields, the vector-type perturbations are therefore irrelevant in the inflationary scenario. The equation for the tensor-type, or gravitational wave perturbations follows from the $`G_j^i`$ component of (7): $$C^{(t)\prime \prime }+2C^{(t)}+k^2C^{(t)}=\kappa _0^2a^2\pi ^{(t)}.$$ (111) For scalar fields we have $`\pi ^{(t)}=0`$. We can recast Eq. (111) as $$v_t^{\prime \prime }+\left(k^2\frac{a^{\prime \prime }}{a}\right)v_t=0,v_t=aC^{(t)},$$ (112) which is of a form similar to Eq. (85). The solution in the large-scale limit is then obtained by ignoring the $`k^2`$ dependent term: $$C^{(t)}(\eta ,𝒌)=A_k+B_k^\eta \frac{1}{a^2}𝑑\eta ,$$ (113) where $`A_k`$ and $`B_k`$ are integration constants. Observe that $`\mathrm{\Psi }^{(v)}`$, $`v^{(v)}`$, $`\pi ^{(v)}`$ and $`C^{(t)}`$ appearing in these equations are gauge-invariant, see Eq. (37). It is interesting to note that we can derive the equations for vector and tensor perturbations without taking into account the scalar fields. Therefore the presence of scalar fields do not formally affect these perturbations. ## VI Conclusion In this paper we presented a general framework for analyzing linear cosmological perturbations in the multicomponent inflation scenario using the gauge-ready approach. Our model consists of multiple scalar fields induced with a positive-definite, but a general field metric, coupled non-minimally to Einstein gravity. The space-time metric is chosen as the perturbed FRW world model. We gave the complete set of perturbation equations in the gauge-ready form, and derived a set of equations for gauge-invariant perturbed order variables. We wrote the equations governing scalar perturbations using generalized forms of slow-roll variables. We then applied canonical quantization to the scalar perturbations and obtained the solutions to the first order in slow-roll. We found that the asymptotic solutions in the super-horizon region are given in terms of Hankel functions. This is similar to the single-field inflation case, except that the order of the Hankel functions is matrix valued. There are a number of possible extensions to our work. First, the immediate next step would be to calculate the power-spectra and spectral indices in realistic models of multiple-field inflation. Second, it would be interesting to interface our approach with the CMBFAST cmbfast or the CAMB camb computer codes and compare with the WMAP results. Third, the methods in this paper may be extended to include generalized gravity as well as hydrodynamical fluids. Fourth, a systematic investigation of adiabatic and isocurvature perturbations in the context of multicomponent inflation can be made using the gauge-ready approach. These, and other issues, will be presented in a future work.
warning/0506/cond-mat0506014.html	ar5iv	text	# Stability of thermodynamic and dynamical order in a system of globally coupled rotors ## 1 Introduction The system of sinusoidally coupled oscillators serves as a prototype model describing various oscillatory phenomena in nature. When the coupling is short-ranged, i.e., between nearest neighbors, the oscillator system describes an array of Josephson junctions, which has been a subject of extensive studies . On the other hand, there are also many systems with long-range couplings in physics and biology. Physiological rhythmic processes may be examples of the latter, which may be modelled as a system of coupled oscillators with the range of coupling being varied, where phase synchronization of the system is an important issue to be understood . Physical examples are diverse, ranging from self-gravitating and plasma systems, where the long-range nature of the gravitational or Coulombic interaction gives rise to difficulty in understanding the systems. A system of globally coupled rotors has thus been proposed and studied to simulate those systems . Here the interaction range is infinite, with the strength scaled with the system size, making the system of the mean-field character and amenable to analytical treatment. In spite of the mean-field nature, however, the system has turned out to exhibit rich features in dynamical and statistical properties. In the canonical ensemble one can find an analytic solution and the system with the ferromagnetic interaction undergoes an equilibrium phase transition at a finite critical temperature, whereas there is no phase transition for the antiferromagnetic interaction. On the other hand, direct simulations in the microcanonical ensemble reveal some interesting features with remarkable differences with the nature of the interaction. Specifically, for the ferromagnetic interaction, the system displays extremely slow relaxation towards the thermodynamic equilibrium. This slow relaxation, dubbed quasi-stationarity, does not coincide with predictions in the canonical ensemble, and thus suggestion has been made that there may exist inequivalence between canonical and microcanonical ensembles. Such quasi-stationarity is observed to survive well below the equilibrium critical temperature and hence has attracted much attention , together with some controversy . In the regime showing quasi-stationarity it has also been reported that the system exhibits aging effects and glassy behavior . For the antiferromagnetic interaction the system exhibits a different type of coherent motion at low temperatures, again only in the microcanonical ensemble : The rotors move in two groups, called the bi-cluster, for a long time, which is explained in terms of the statistical equilibrium of the effective Hamiltonian obtained after averaging out fast variables. In a recent work we have employed a novel approach that treats the system in a unified framework of microcanonical and canonical ensembles . Starting from the set of Langevin equations describing dissipative dynamics of a system (canonical ensemble) and the corresponding Fokker-Planck equation (FPE), we have pointed out that the nondissipative Hamiltonian dynamics (microcanonical ensemble) may be described as a limiting case of the vanishing damping coefficient. Thereupon we have been able to find a class of solution for the incoherent phase depending on the ensemble, some of which are neutrally stable even below the equilibrium critical temperature. This neutral stability has then been suggested to be a plausible physical explanation as to the origin of the quasi-stationarity observed in numerical experiments. In this paper we further extend the stability analysis of the previous work to the ferromagnetic coherent phase (with thermodynamic order) and to the system with the antiferromagnetic interaction. For the latter, we attempt to provide an alternative view of the bi-cluster phase observed in the antiferromagnetic system, as dynamical order allowed by the rotating solution of the FPE. This rotating solution is found to be neutrally stable down to zero temperature. Furthermore, the rotating solution can give rise to any degree of clustering, if the initial condition is appropriately chosen, in addition to bi-clustering. It would thus be of interest to probe such multi-cluster motions as tri-clustering, etc., by means of numerical simulations. This paper is organized as follows: In Sec. 2 we describe how the system of globally coupled rotors can be treated in a unified framework from the set of Langevin equations and the corresponding FPE. Various solutions of the FPE are given in Sec. 3. It is shown that multi-cluster solutions emerge, manifesting dynamical order for the non-stationary rotating solution of the FPE. Section 4 is devoted to the stability analysis of the stationary solutions, with emphasis on the ferromagnetically coherent phase (with single cluster motion or thermodynamic order). The stability analysis of the non-stationary solution is presented in Sec. 5, with a special focus on the antiferromagnetic case. Finally, a brief summary is given in Sec. 6. ## 2 System of Coupled Rotors We consider a system of $`N`$ classical rotors, each of which is described by its phase angle and coupled sinusoidally to others. The dynamics of the coupled rotor system is governed by the set of equations of motion for the phase $`\varphi _i`$ ($`i=1,\mathrm{},N`$) of the $`i`$th rotor: $$M\ddot{\varphi _i}+\underset{j}{}J_{ij}\mathrm{sin}(\varphi _i\varphi _j)=0,$$ (1) where $`M`$ is the inertia of each rotor and $`J_{ij}`$ represents the coupling strength between rotors $`i`$ and $`j`$. With the introduction of the canonical momentum $`p_i=M\dot{\varphi _i}`$, the above equations are transformed into a set of canonical equations: $$\dot{\varphi _i}=\frac{_N}{p_i},\dot{p_i}=\frac{_N}{\varphi _i}$$ (2) with the $`N`$-particle Hamiltonian $$_N=\underset{i}{}\frac{p_i^2}{2M}\underset{i<j}{}J_{ij}\mathrm{cos}(\varphi _i\varphi _j),$$ (3) on which the microcanonical description is based. On the other hand, in the canonical description the system is in contact with a heat reservoir of temperature $`T`$ and described, in a most general way, by the set of Langevin equations: $$M\ddot{\varphi _i}+\mathrm{\Gamma }\dot{\varphi _i}+\underset{j}{}J_{ij}\mathrm{sin}(\varphi _i\varphi _j)=\eta _i,$$ (4) where $`\mathrm{\Gamma }`$ is the damping coefficient and the Gaussian white noise $`\eta _i(t)`$ is characterized by the average $`\eta _i(t)=0`$ and the correlation $`\eta _i(t)\eta _j(t^{})=2\mathrm{\Gamma }T\delta _{ij}\delta (tt^{})`$. To derive the corresponding FPE, we write the equations of motion in the form $`\dot{\varphi _i}`$ $`=`$ $`{\displaystyle \frac{p_i}{M}}`$ $`\dot{p_i}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }}{M}}p_i{\displaystyle \underset{j}{}}J_{ij}\mathrm{sin}(\varphi _i\varphi _j)+\eta _i.`$ (5) It is then straightforward to derive, via the standard procedure , the FPE for the probability distribution $`P(\varphi _i,p_i,t)`$: $`{\displaystyle \frac{P}{t}}`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{p_i}{M}}{\displaystyle \frac{P}{\varphi _i}}+{\displaystyle \underset{i}{}}{\displaystyle \frac{}{p_i}}`$ (6) $`\times `$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }}{M}}p_i+{\displaystyle \underset{j}{}}J_{ij}\mathrm{sin}(\varphi _i\varphi _j)+\mathrm{\Gamma }T{\displaystyle \frac{}{p_i}}\right]P.`$ One may also derive the FPE for the Hamiltonian dynamics, which just reads Eq. (6) with $`\mathrm{\Gamma }=0`$. While reflecting that Eq. (4) with $`\mathrm{\Gamma }`$ set equal to zero reduces to Eq. (1), this suggests that Eq. (6) should provide the starting point for both descriptions: the microcanonical one ($`\mathrm{\Gamma }=0`$) and the canonical one ($`\mathrm{\Gamma }0`$). In particular, the stationary solution of Eq. (6) is given by the canonical distribution $`P^{(0)}(\varphi _i,p_i)e^{_N/T}`$, describing equilibrium, with the very Hamiltonian in Eq. (3) regardless of $`\mathrm{\Gamma }`$ being zero or not. Note, however, that unlike the canonical ensemble where $`T`$ represents the given temperature, in the microcanonical ensemble $`T`$ still remains as an arbitrary parameter. In the latter, one may adjust $`T`$ to the average kinetic energy, which allows the interpretation of $`T`$ as the temperature. This prescription thus establishes correspondence between the two ensembles. Note also that in the zero-temperature limit ($`T0`$), Eq. (4) reduces to the Caldirola-Kanai Hamiltonian dynamics , which needs external driving to have a nontrivial stationary state. In order to measure a variety of coherence in the system, we conveniently introduce the generalized order parameter $`\mathrm{\Delta }^{(\mathrm{})}`$ defined by $$\frac{1}{N}\underset{i}{\overset{N}{}}e^{i\mathrm{}\varphi _i}\mathrm{\Delta }^{(\mathrm{})}e^{i\theta _{\mathrm{}}}.$$ (7) Apart from the global phase $`\theta _{\mathrm{}}`$, non-vanishing values of the order parameter $`\mathrm{\Delta }^{(\mathrm{})}`$ imply that rotors move as clusters, since rotors separated with phase angle $`2\pi /\mathrm{}`$ make contributions to $`\mathrm{\Delta }^{(\mathrm{})}`$. It thus can be used as a measure of the distribution of rotors, particularly, the degree of clustering. For instance, a non-vanishing value for $`\mathrm{}=1`$ corresponds to the emergence of a mono-cluster (or magnetization), that for $`\mathrm{}=2`$ corresponds to bi-cluster formation (with separation of $`\pi `$), and so on. Note that the $`\mathrm{}=2`$ case may be regarded as the analogue of staggered magnetization in the short-ranged model. ## 3 Stationary and Non-Stationary Solutions In the infinite-range limit ($`J_{ij}=J/N`$ with $`N\mathrm{}`$), we use Eq. (7) for $`\mathrm{}=1`$, $$\mathrm{\Delta }^{(1)}=\frac{1}{N}\underset{i}{}e^{i(\varphi _i\theta _1)},$$ (8) and decouple the set of the equations of motion into a single-particle equation $$M\ddot{\varphi _i}+\mathrm{\Gamma }\dot{\varphi _i}+J\mathrm{\Delta }^{(1)}\mathrm{sin}(\varphi _i\theta _1)=\eta _i,$$ (9) satisfied by all rotors. Henceforth we therefore drop the rotor index $`i`$ in Eq. (9), which leads to the standard FPE for the single-rotor probability distribution $`P(\varphi ,p,t)`$: $`{\displaystyle \frac{P}{t}}=`$ $``$ $`{\displaystyle \frac{p}{M}}{\displaystyle \frac{P}{\varphi }}+J\mathrm{\Delta }^{(1)}\mathrm{sin}(\varphi \theta _1){\displaystyle \frac{P}{p}}`$ (10) $`+`$ $`\mathrm{\Gamma }{\displaystyle \frac{}{p}}\left[{\displaystyle \frac{p}{M}}+T{\displaystyle \frac{}{p}}\right]P.`$ In the absence of damping ($`\mathrm{\Gamma }=0`$), this reduces to the FPE for the microcanonical ensemble: $$\frac{P}{t}=\frac{p}{M}\frac{P}{\varphi }+J\mathrm{\Delta }^{(1)}\mathrm{sin}(\varphi \theta _1)\frac{P}{p},$$ (11) which is also referred to as the Vlasov equation in some literature . In terms of this probability distribution, the generalized order parameter is defined to be $$\mathrm{\Delta }^{(\mathrm{})}e^{i\theta _{\mathrm{}}}=e^{i\mathrm{}\varphi }=𝑑p𝑑\varphi e^{i\mathrm{}\varphi }P(\varphi ,p,t).$$ (12) ### 3.1 Stationary Solutions For the sake of completeness, we briefly review the results for stationary solutions of the FPE . As pointed out for the general case, both Eqs. (10) and (11) support the same stationary ($`P/t=0`$) solution: $$P^{(0)}(\varphi ,p)=\frac{1}{𝒵}e^{/T}$$ (13) with the single-particle Hamiltonian $$=\frac{p^2}{2M}J\mathrm{\Delta }^{(1)}\mathrm{cos}(\varphi \theta _1),$$ (14) where the overall phase $`\theta _1`$ manifests the global $`U(1)`$ symmetry. It is thus expected that both ensembles exhibit the same equilibrium behavior. One, however, should recall again that here $`T`$ is given in Eq. (10) (for the canonical ensemble) but remains arbitrary in Eq. (11) (for the microcanonical ensemble). In the microcanonical ensemble the temperature should be defined as a measure of the average kinetic energy according to $`p^2/2MT/2`$. The partition function is determined by normalization: $$𝒵=𝑑p\frac{d\varphi }{2\pi }e^{/T}.$$ (15) For later use, we first describe some equilibrium properties of the globally coupled rotors through the use of the single-particle model. Defining $`xJ\mathrm{\Delta }^{(1)}/T`$ and making use of the expansion $$e^{x\mathrm{cos}(\varphi \theta _1)}=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}I_n(x)e^{in(\varphi \theta _1)}$$ (16) with $`I_n(x)`$ being the modified Bessel function of the $`n`$-th order, we evaluate the partition function as $$𝒵=\sqrt{2\pi MT}I_0(x).$$ (17) We emphasize again that this approach based on the FPE provides a unified description of microcanonical and canonical ensembles and both ensembles generate the same equilibrium behavior, determined by the same distribution $`P^{(0)}(\varphi ,p)`$. Namely, in both ensembles the generalized order parameter in equilibrium is given by $`\mathrm{\Delta }^{(\mathrm{})}e^{i\theta _{\mathrm{}}}`$ $`=`$ $`e^{i\mathrm{}\varphi }={\displaystyle 𝑑p\frac{d\varphi }{2\pi }P^{(0)}(\varphi ,p)e^{i\mathrm{}\varphi }}.`$ (18) With the expansion in Eq. (16) and integration over $`\varphi `$, the order parameter reads $$\mathrm{\Delta }^{(\mathrm{})}=\frac{I_{\mathrm{}}(x)}{I_0(x)}.$$ (19) Note here that $`\mathrm{\Delta }^{(\mathrm{})}`$ has an explicit dependence on the coherence order parameter $`\mathrm{\Delta }^{(1)}`$ through $`xJ\mathrm{\Delta }^{(1)}/T`$. For $`\mathrm{}=1`$, describing the emergence of coherence (the mono-cluster as thermodynamic order), Eq. (18) becomes an equation to be solved self-consistently: $$\frac{T}{J}x=\frac{I_1(x)}{I_0(x)}.$$ (20) This self-consistency equation determines whether the system exhibits coherence: The ordered phase ($`\mathrm{\Delta }^{(1)}0`$) emerges when $`T/J`$ is smaller than the slope of $`I_1(x)/I_0(x)`$ at $`x=0`$, which is 1/2. Accordingly, the ferromagnetic system ($`J>0`$) undergoes a phase transition at the critical temperature $`T_c=J/2`$. In the case of antiferromagnetic coupling ($`J<0`$), on the other hand, Eq. (20) becomes $`Tx/\|J\|=I_1(x)/I_0(x)`$, leading to the only solution $`x=0`$. It is thus concluded that the antiferromagnetic system has no phase transition at finite temperatures (no mono-cluster). It is obvious in Eq. (19) that $`\mathrm{\Delta }^{(\mathrm{})}`$ for higher values of $`\mathrm{}`$ can assume nonzero values only for $`\mathrm{\Delta }^{(1)}0`$; this implies that only the ferromagnetic system can develop all degrees of clustering below $`T_c`$. This is not surprising since the mono-cluster phase has a $`2\pi `$ symmetry and therefore invariant under any rotations of multiples of $`2\pi `$, which in turn gives rise to nonzero $`\mathrm{\Delta }^{(\mathrm{})}`$. For $`\mathrm{\Delta }^{(1)}=0`$, describing the incoherent phase, the single-particle Hamiltonian (14) has only the kinetic energy term, thus reducing the canonical distribution $`P^{(0)}(\varphi ,p)`$ to the Maxwell distribution for both ensembles. Unlike Eq. (10), however, Eq. (11), the FPE in the microcanonical ensemble, allows an extra solution of the form $`P^{(0)}(\varphi ,p)=f_0(p)`$, an arbitrary function of $`p`$ without $`\varphi `$-dependence, including the Maxwell distribution as a special case . The only constraint is the normalization, and the distribution $`P^{(0)}(\varphi ,p)`$ uniform in $`\varphi `$ guarantees $`\mathrm{\Delta }^{(1)}=0`$. As a result, $`\mathrm{\Delta }^{(\mathrm{})}`$ vanishes for all values of $`\mathrm{}`$ as well, and no multi-clustering is allowed by this type of stationary solution present only in the microcanonical ensemble. ### 3.2 Rotating Solutions In addition to the stationary solutions presented above, the FPE in the microcanonical ensemble also carries non-stationary solutions which have some significance for the antiferromagnetic system. For $`\mathrm{\Delta }^{(1)}=0`$, Eq. (11) becomes $$\frac{P}{t}=\frac{p}{M}\frac{P}{\varphi },$$ (21) which has a solution of the general form $`P^{(0)}(\varphi ,p,t)=u(\varphi pt/M,p)`$. This is a rotating solution in the sense that the phase grows continuously with time with a continuous frequency spectrum ($`\omega p/M`$). Requiring periodicity in $`\varphi `$, we write $$P^{(0)}(\varphi ,p,t)=\underset{k}{}e^{ik(\varphi \frac{p}{M}t)}F_k(p),$$ (22) where $`F_k(p)`$ is an arbitrary function of $`p`$ satisfying $`F_{\pm 1}(p)=0`$ due to the condition $`\mathrm{\Delta }^{(1)}=0`$. The generalized order parameter for this solutions is computed according to $`\mathrm{\Delta }^{(\mathrm{})}e^{i\theta _{\mathrm{}}}`$ $`=`$ $`{\displaystyle 𝑑p\frac{d\varphi }{2\pi }\underset{k}{}e^{i\mathrm{}\varphi }e^{ik(\varphi \frac{p}{M}t)}F_k(p)}`$ (23) $`=`$ $`{\displaystyle 𝑑pe^{+i\mathrm{}\frac{p}{M}t}F_{\mathrm{}}(p)},`$ which shows that the higher-order moment $`\mathrm{\Delta }^{(\mathrm{})}`$ in general does not vanish unless $`F_{\mathrm{}}(p)=0`$. Thus far there is no difference between the ferromagnetic and the antiferromagnetic couplings. As will be shown later, however, the stability of the rotating solution differs substantially, depending on the nature of the interaction. The rotating solution exists only for $`\mathrm{\Delta }^{(1)}=0`$, regardless of whether the system is in equilibrium or not. While such an incoherent phase appears only at high temperatures in the ferromagnetic case, $`\mathrm{\Delta }^{(1)}`$ remains always zero at all temperature ranges in the antiferromagnetic case. Moreover, the rotation frequency gets higher as the order of the moment increases. This suggests that at low temperatures where thermal fluctuations are small, the phases with non-vanishing higher moments (high degrees of clustering) are easier to observe in the antiferromagnetic system than in the ferromagnetic one. In fact, this is precisely what has been seen in recent numerical simulations, which reported the bi-cluster phase in the antiferromagnetic system at very low temperatures . The bi-cluster state, with two clusters separated by angle $`\pi `$, may be obtained with suitable choices of $`F_k(p)`$ in Eq. (22). For example, with the choice $`F_{2k}(p)=F(p)`$ and $`F_{2k+1}(p)=0`$, we obtain $$P^{(0)}(\varphi ,p,t)=\pi F(p)\left[\delta \left(\varphi \frac{p}{M}t\right)+\delta \left(\varphi \frac{p}{M}t+\pi \right)\right].$$ (24) For another choice, say $`F_{2k}(p)=(1)^kF(p)`$ and $`F_{2k+1}(p)=0`$, one obtains $$P^{(0)}(\varphi ,p,t)=\pi F(p)\left[\delta \left(\varphi \frac{p}{M}t+\frac{\pi }{2}\right)+\delta \left(\varphi \frac{p}{M}t\frac{\pi }{2}\right)\right].$$ (25) All these are shown to be neutrally stable in Sec. 5. Equation (23) further indicates that there can exist higher-order multi-cluster phases (for $`\mathrm{}=2,3,\mathrm{}`$) as well, if appropriate choices for $`F_k(p)`$ are made. Recall again that the multi-cluster phase does not occur for stationary solutions since $`\mathrm{\Delta }^{(\mathrm{})}=0`$ for time-independent solutions. In other words, the multi-cluster must rotate with the frequency higher as the number of clusters grows; this suggests that the multi-cluster with large $`\mathrm{}`$ should be difficult to observe. ## 4 Stability of Stationary States In the previous work , we have already shown that the stability of the incoherent phase depends on the solutions of the FPE, providing a plausible explanation as to the physical origin of the quasi-stationarity. We now extend the analysis further to include the case of the coherent phase. For this purpose, we write the FPE, setting $`\mathrm{\Delta }^{(1)}\mathrm{\Delta }`$ and $`\theta _1\theta `$, in the form $$\frac{P}{t}=\frac{p}{M}\frac{P}{\varphi }+J\mathrm{\Delta }\mathrm{sin}(\varphi \theta )\frac{P}{p}+\mathrm{\Gamma }\frac{}{p}\left[\frac{p}{M}+T\frac{}{p}\right]P,$$ (26) To probe the stability, we add a small perturbation to write $$P(\varphi ,p,t)=P_0(\varphi ,p,t)+f(\varphi ,p,t)$$ (27) and accordingly $`\mathrm{\Delta }(t)`$ $`=`$ $`\mathrm{\Delta }_0(t)+\mathrm{\Delta }_1(t)`$ (28) $`=`$ $`{\displaystyle 𝑑p\frac{d\varphi }{2\pi }e^{i(\varphi \theta )}\left[P_0(\varphi ,p,t)+f(\varphi ,p,t)\right]}.`$ Substituting these into (26), one obtains, to the lowest order, $`{\displaystyle \frac{f}{t}}=`$ $``$ $`{\displaystyle \frac{p}{M}}{\displaystyle \frac{f}{\varphi }}+J\mathrm{\Delta }_1\mathrm{sin}(\varphi \theta ){\displaystyle \frac{P_0}{p}}`$ (29) $`+`$ $`J\mathrm{\Delta }_0\mathrm{sin}(\varphi \theta ){\displaystyle \frac{f}{p}}+\mathrm{\Gamma }{\displaystyle \frac{}{p}}\left({\displaystyle \frac{p}{M}}+T{\displaystyle \frac{}{p}}\right)f.`$ Since $`f(\varphi ,p,t)`$ and $`\mathrm{\Delta }_1(t)`$ are periodic in $`\varphi `$, one can Fourier expand them in plane waves: $$f(\varphi ,p,t)=\underset{k}{}𝑑\omega e^{i(k\varphi \omega t)}\stackrel{~}{f}_k(p,\omega )$$ (30) and $`\mathrm{\Delta }_1(t)`$ $`=`$ $`{\displaystyle 𝑑p\frac{d\varphi }{2\pi }e^{i(\varphi \theta )}f(\varphi ,p,t)}`$ (31) $`=`$ $`{\displaystyle 𝑑\omega e^{i\omega t}𝑑p\stackrel{~}{f}_1(p,\omega )},`$ where the integration over $`\varphi `$ has been performed. Note here that the perturbed order parameter is proportional only to $`\stackrel{~}{f}_1(p,\omega )`$ (or to $`\stackrel{~}{f}_{+1}(p,\omega )`$ if the order parameter has been defined to be $`\mathrm{\Delta }=e^{i(\varphi \theta )}`$). Inserting these expressions into Eq. (29) and collecting coefficients of $`e^{i(k\varphi \omega t)}`$, one finds the relations satisfied by the Fourier coefficients $`\stackrel{~}{f}_k(p,\omega )`$. In the case of ferromagnetic coupling, the coherent phase ($`\mathrm{\Delta }_00`$) arises at temperatures below $`T_c`$, regardless of the presence of damping. The stationary solution in Eq. (13) can be written, with the help of Eqs. (16), (17), and (19), in the form $`P_0(\varphi ,p)`$ $`=`$ $`f_M(p){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{I_n(x)}{I_0(x)}}e^{in(\varphi \theta )}`$ (32) $`=`$ $`f_M(p){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\mathrm{\Delta }^{(n)}(x)e^{in(\varphi \theta )},`$ where $`f_M(p)(2\pi MT)^{1/2}\mathrm{exp}(p^2/2MT)`$ is the Maxwell distribution and $`xJ\mathrm{\Delta }_0/T`$ as before. When $`x=0`$, the above equation simply reduces to the Maxwell distribution, which is stable at temperatures above $`T_c`$. Our concern now is how the coherent phase gets its stability as the temperature is lowered below the critical temperature. Putting Eqs. (31) and (32) into Eq. (29), one obtains the following equation for the Fourier coefficients $`\stackrel{~}{f}_k(p,\omega )`$: $`\left(\omega {\displaystyle \frac{kp}{M}}\right)\stackrel{~}{f_k}`$ $``$ $`{\displaystyle \frac{J\mathrm{\Delta }_0}{2}}{\displaystyle \frac{}{p}}(\stackrel{~}{f}_{k1}\stackrel{~}{f}_{k+1})i\mathrm{\Gamma }{\displaystyle \frac{}{p}}\left({\displaystyle \frac{p}{M}}+T{\displaystyle \frac{}{p}}\right)\stackrel{~}{f}_k`$ (33) $`=`$ $`{\displaystyle \frac{J}{2}}[\mathrm{\Delta }^{(k1)}(x)\mathrm{\Delta }^{(k+1)}(x)]f_M^{}(p){\displaystyle 𝑑p^{}\stackrel{~}{f}_1}.`$ We note here that the emergence of coherence contributes to the off-diagonal term in Eq. (33) and to the appearance of higher-order generalized order parameters, making the stability analysis non-trivial. When $`\omega kp/M=0`$, we have a continuous spectrum, and for $`\mathrm{\Gamma }=0`$, Eq. (33) becomes $$\frac{J\mathrm{\Delta }_0}{2}\frac{}{p}\stackrel{~}{f}_{k1}=\frac{J}{2}\mathrm{\Delta }^{(k1)}(x)f_M^{}(p)\left[𝑑p^{}\stackrel{~}{f}_1\right].$$ (34) It is easy to show by direct substitution that this equation has a solution of the form $$\stackrel{~}{f}_k(p,\omega )=\{\begin{array}{cc}f_M(p)h_k(\omega ),\hfill & \text{for }k\pm 1,0,2\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$ (35) It is of interest to note this is also the solution for $`\mathrm{\Gamma }0`$ as well, since the term including $`\mathrm{\Gamma }`$ vanishes for the Maxwell distribution. For $`\omega kp/M0`$, we have a discrete spectrum and may not solve the equation for the general case. Still we may proceed further if we take the phase-only perturbation, namely, $`f(p,\varphi ,t)=f_M(p)h(\varphi ,t)`$. We then have the Fourier coefficient $`\stackrel{~}{f}_k(p,\omega )=f_M(p)h_k(\omega )`$ with $`h_k(\omega )`$ being the Fourier coefficient of $`h(\varphi ,t)`$, which in turn gives $`(/p)\stackrel{~}{f}_k(p,\omega )=f_M^{}(p)h_k(\omega )`$ and $`𝑑p\stackrel{~}{f}_k(p,\omega )=h_k(\omega )`$. Further, the term including $`\mathrm{\Gamma }`$ vanishes identically in this case. Dividing Eq. (33) by $`\omega kp/M`$ and integrating over $`p`$, we obtain $`h_k(\omega )=`$ $``$ $`[\mathrm{\Delta }^{(k1)}(x)\mathrm{\Delta }^{(k+1)}(x)]\chi _k(\omega )\stackrel{~}{h}_1(\omega )`$ (36) $`+`$ $`2\mathrm{\Delta }_0\chi _k(\omega )[h_{k+1}(\omega )h_{k1}(\omega )],`$ where we have introduced the $`k`$-dependent response function $$\chi _k(\omega )=\frac{J}{2}𝑑p\frac{f_M^{}(p)}{\omega +kp/M}$$ (37) and used Eq. (75). Some properties of this response function, which is frequency-dependent, are discussed separately in the Appendix. For $`x0`$, the recursion relation for the modified Bessel functions : $$I_{k1}(x)I_{k+1}(x)=\frac{2k}{x}I_k(x)$$ leads Eq. (36) to take the form $$h_k2\mathrm{\Delta }_0\chi _k(\omega )(h_{k+1}h_{k1})=\frac{2k}{x}\mathrm{\Delta }^{(k)}(x)\chi _k(\omega )h_1,$$ (38) which needs to be solved. For $`k=0`$, from Eqs. (36) and (73), we find $`h_0=0`$, implying the absence of a constant term in the perturbation. Noting $`\mathrm{\Delta }^{(k)}(x)=\mathrm{\Delta }^{(k)}(x)`$ and $`\chi _k(\omega )=\chi _k(\omega )`$, we write the difference equation in the matrix form: $$\mathrm{\Lambda }\left(\begin{array}{c}h_1\\ h_2\\ h_3\\ \mathrm{}\end{array}\right)=0$$ (39) with the matrix $$\mathrm{\Lambda }=\left(\begin{array}{cccc}1+(2/x)\mathrm{\Delta }^{(1)}\chi _1& \mathrm{\Delta }_0\chi _1& 0& \mathrm{}\\ \mathrm{\Delta }_0\chi _2+(4/x)\mathrm{\Delta }^{(2)}\chi _2& 1& \mathrm{\Delta }_0\chi _2& \mathrm{}\\ (6/x)\mathrm{\Delta }^{(3)}\chi _3& \mathrm{\Delta }_0\chi _3& 1& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$ (40) Here we have included the terms with only negative $`k`$ values, reflecting that all order parameters are defined by Eq. (18). In order to have non-trivial solutions for $`\stackrel{}{h}=(h_1,h_2,h_3,\mathrm{})`$, one should have the vanishing determinant: $$\epsilon (\omega )\text{det}\mathrm{\Lambda }=0.$$ (41) Let us first consider the limit $`\mathrm{\Delta }_00`$ or $`x=J\mathrm{\Delta }_0/T0`$, corresponding to the incoherent phase. In this limit all the off-diagonal terms vanish, since $`I_0(x)1`$ and $`I_n(x)(x/2)^n`$ so that $`\mathrm{\Delta }^{(n)}(x/2)^n`$. Equation (39) then becomes $$\left(\begin{array}{cccc}1+\chi _1& 0& 0& \mathrm{}\\ 0& 1& 0& \mathrm{}\\ 0& 0& 1& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)\left(\begin{array}{c}h_1\\ h_2\\ h_3\\ \mathrm{}\end{array}\right)=0,$$ (42) which leads to $$1+\chi _1(\omega )=1+\chi (\omega )1+\frac{JM}{2}\stackrel{~}{\chi }(\omega )=0$$ (43) for non-vanishing $`h_1`$, while all other $`h`$’s are zero. The detailed analytic properties of the reduced response function $`\stackrel{~}{\chi }_(\omega )(2/JM)\chi _(\omega )`$, with the complex frequency $`\omega =\omega _r+i\omega _i`$, are presented in the Appendix. Equation (43) describes the condition for the incoherent phase with the Maxwell distribution, which is stable/unstable above/below $`T_c`$ . For $`x0`$, in principle we have to solve Eq. (41) including all the terms in Eq. (40). Since this is very formidable, we instead consider just a few terms to explore how the stability of the solution changes. To this end, we write $`\epsilon (\omega )\epsilon ^{(m)}(\omega )`$, the determinant obtained when the first $`m`$ Fourier components are kept. With only the first Fourier component $`h_1`$ considered, Eq. (41) obtains the form $`\epsilon ^{(1)}(\omega )`$ $`=`$ $`1+{\displaystyle \frac{2}{x}}\mathrm{\Delta }^{(1)}\chi _1(\omega )`$ (44) $`=`$ $`1+{\displaystyle \frac{2T}{J}}\chi _(\omega )`$ $`=`$ $`1+TM\stackrel{~}{\chi }(\omega )=0.`$ for which Eqs. (91) to (93) yield $`\omega _i=0`$ as the only solution. Comparison of Eq. (44) with Eq. (43) shows the correspondence $`T=J/2=T_c`$; this indicates that the solution is neutrally stable at the critical temperature, below which coherence develops. Including the next component $`h_2`$, one has $$\epsilon ^{(2)}(\omega )=\epsilon ^{(1)}(\omega )+\mathrm{\Delta }_0(\mathrm{\Delta }_0+\frac{4}{x}\mathrm{\Delta }^{(2)})\chi _1(\omega )\chi _2(\omega )=0,$$ (45) which, with $`\mathrm{\Delta }^{(1)}=\mathrm{\Delta }_0`$ and $`\chi _k(\omega )=\chi (\omega /k)/k`$, becomes $$\epsilon ^{(2)}(\omega )=1+T\stackrel{~}{\chi }(\omega )+\frac{T^2}{8}\left[x^2+4x\frac{I_2(x)}{I_1(x)}\right]\stackrel{~}{\chi }(\omega )\stackrel{~}{\chi }(\omega /2)=0.$$ (46) Since $`\stackrel{~}{\chi }(\omega )`$ and $`\stackrel{~}{\chi }(\omega /2)`$ have the same pole structure, the real and the imaginary parts of Eq. (46) read $`\text{Re}\epsilon ^{(2)}(\omega )`$ $``$ $`1+T\text{Re}\stackrel{~}{\chi }(\omega )+{\displaystyle \frac{T^2}{8}}\left[x^2+4x{\displaystyle \frac{I_2(x)}{I_1(x)}}\right]`$ (47) $`\times [\text{Re}\stackrel{~}{\chi }(\omega )\text{Re}\stackrel{~}{\chi }(\omega /2)\text{Im}\stackrel{~}{\chi }(\omega )\text{Im}\stackrel{~}{\chi }(\omega /2)]=0`$ $`\text{Im}\epsilon ^{(2)}(\omega )`$ $``$ $`T\text{Im}\stackrel{~}{\chi }(\omega )+{\displaystyle \frac{T^2}{8}}\left[x^2+4x{\displaystyle \frac{I_2(x)}{I_1(x)}}\right]`$ (48) $`\times [\text{Im}\stackrel{~}{\chi }(\omega )\text{Re}\stackrel{~}{\chi }(\omega /2)+\text{Re}\stackrel{~}{\chi }(\omega )\text{Im}\stackrel{~}{\chi }(\omega /2)]=0.`$ In the Appendix it is shown that $`\omega _r=0`$ is a solution of $`\text{Im}\stackrel{~}{\chi }(\omega )=0`$, implying that this is also a solution of $`\text{Im}\epsilon ^{(2)}(\omega )=0`$. For $`\omega _i>0`$, Eq. (47) becomes $$f(y)1=\frac{1}{8}\left[x^2+4x\frac{I_2(x)}{I_1(x)}\right]f(y)f(y/2)$$ (49) with $`y\omega _i\sqrt{M/2T}`$. As $`y`$ increases from zero to arbitrarily large values, the left-hand side of Eq. (49) decreases monotonically from zero to $`1`$ while the right-hand side is positive-definite for $`y>0`$. This suggests that there is no solution for $`\omega _i>0`$ to make the system unstable. For $`\omega _i=0`$, which corresponds to the neutral stability, we have $`\text{Re}\stackrel{~}{\chi }(\omega )=1/T`$ and thus Eq. (47) reads $$x^2+4x\frac{I_2(x)}{I_1(x)}=0,$$ (50) which leads to $`x=0`$ for the critical case. For $`\omega _i=\|\omega _i\|<0`$, for which the system becomes stable, Eq. (47) takes the form $$g(y)1=\frac{1}{8}\left[x^2+4x\frac{I_2(x)}{I_1(x)}\right]g(y)g(y/2).$$ (51) As shown in the Appendix, $`g(y)`$ is a monotonically increasing function of $`y`$ from unity to infinity, and accordingly, $`g(y)g(y/2)`$ is also a monotonically increasing function of $`y`$ in the same domain. Since the left-hand side of the above equation is monotonically increasing from zero to arbitrarily large values, Eq. (51) allows a solution only for some range of $`x`$ values. We have determined numerically the range of $`x`$ values, in which there exits a solution for $`y>0`$, to find $$x\frac{J\mathrm{\Delta }_0}{T}<x_c^{(2)}1.32.$$ (52) This indicates that the coherent solution is stable only at temperatures above $`T_0`$, at which $`\mathrm{\Delta }_0(T)`$ and $`x_cT/J`$ meet. We see that the stable region does not extend to the zero temperature, presumably because we have included only the second component in our analysis (the first component is trivial). This may be resolved if one include higher components. Adding the third component $`h_3`$ leads to the following equation $`\epsilon ^{(3)}(\omega )`$ $`=`$ $`\epsilon ^{(2)}(\omega )+{\displaystyle \frac{1}{24}}J^2\mathrm{\Delta }_0^2\left[1+T\stackrel{~}{\chi }(\omega )+3T{\displaystyle \frac{I_3(x)}{I_1(x)}}\stackrel{~}{\chi }(\omega )\right]\stackrel{~}{\chi }(\omega /2)\stackrel{~}{\chi }(\omega /3)`$ (53) $`=`$ $`0`$ from which one can perform the similar analysis to find $`\left[1+{\displaystyle \frac{x^2}{24}}f(y/2)f(y/3)\right][f(y)1]`$ $`={\displaystyle \frac{1}{8}}\left[x^2+4x{\displaystyle \frac{I_2(x)}{I_1(x)}}{\displaystyle \frac{I_3(x)}{I_1(x)}}f(y/3)\right]f(y)f(y/2)`$ (54) for $`\omega _i>0`$ and $`\left[1+{\displaystyle \frac{x^2}{24}}g(y/2)g(y/3)\right][g(y)1]`$ $`={\displaystyle \frac{1}{8}}\left[x^2+4x{\displaystyle \frac{I_2(x)}{I_1(x)}}{\displaystyle \frac{I_3(x)}{I_1(x)}}g(y/3)\right]g(y)g(y/2)`$ (55) for $`\omega _i<0`$. Again, Eq. (4) does not have a solution for positive $`y`$ since the left-hand side is less than zero while the right-hand side is greater than zero. Equation (4) is found to have a solution for $`x<x_c^{(3)}1.51`$. Note here that $`x_c`$ is increased substantially once the third component is included, implying that $`T_0`$, above which the coherent solution is stable, is lowered. We have performed this analysis, including up to four components, and confirmed that this trend persists; this suggests the plausible conjecture that the coherent solution is stable down to zero temperature if all the Fourier components are included. We now turn our attention to the stability of the antiferromagnetic system for which there is no equilibrium order ($`\mathrm{\Delta }_0=0`$ or $`x=0`$). The stability equation reads, for $`J=\|J\|`$, $$1\frac{\|J\|}{2}\stackrel{~}{\chi }_(\omega )=0,$$ (56) which, depending on the sign of $`\omega _i`$ (with $`\omega _r=0`$), becomes $$\{\begin{array}{cc}1+(\|J\|/2T)f(y)=0\hfill & \text{for}\omega _i>0\hfill \\ 1+\|J\|/2T=0\hfill & \text{for}\omega _i=0\hfill \\ 1+(\|J\|/2T)g(y)=0\hfill & \text{for}\omega _i<0.\hfill \end{array}$$ (57) None of these equations has a solution, since $`f(y)`$ and $`g(y)`$ is positive-definite. This means that the antiferromagnetic system cannot have self-sustained deviation in the absence of the perturbation with a discrete spectrum. On the other hand, with the continuous spectrum $`\omega =\omega _r=kp/M`$, the system is neutrally stable at all temperatures. ## 5 Stability of Non-Stationary States As mentioned in Section 3, the non-stationary solution exists only in the microcanonical ensemble ($`\mathrm{\Gamma }=0`$) with a continuous frequency spectrum, which can develop spontaneously. Our concern in this section is the stability of this solution, especially in the case of the antiferromagnetic interaction. Equation (29) for stability reads, with $`\mathrm{\Delta }_0=\mathrm{\Gamma }=0`$, $$\frac{f}{t}=\frac{p}{M}\frac{f}{\varphi }+J\mathrm{\Delta }_1\mathrm{sin}\varphi \frac{P_0}{p},$$ (58) where $`P_0`$ is given by Eq. (22). The last term in the above equation obtains the form $`J\mathrm{\Delta }_1\mathrm{sin}\varphi {\displaystyle \frac{P_0}{p}}`$ $`=`$ $`{\displaystyle \frac{J}{2i}}{\displaystyle 𝑑\omega e^{i\omega t}\left[𝑑p^{}\stackrel{~}{f}_1(p^{},\omega )\right](e^{i\varphi }e^{i\varphi })}`$ (59) $`\times {\displaystyle \underset{k}{}}e^{ik\varphi }{\displaystyle \frac{}{p}}\left[\mathrm{exp}(i{\displaystyle \frac{kp}{M}}t)F_k(p)\right]`$ $`=`$ $`{\displaystyle \frac{J}{2i}}{\displaystyle \underset{k}{}}{\displaystyle \frac{}{p}}{\displaystyle 𝑑\omega \left[(e^{i(k+1)\varphi }e^{i(k1)\varphi })\right]\mathrm{exp}\left[i(\omega +\frac{kp}{M})t\right]}`$ $`\times F_k(p){\displaystyle 𝑑p^{}\stackrel{~}{f}_1(p^{},\omega )}`$ $`=`$ $`{\displaystyle \frac{J}{2i}}{\displaystyle \underset{k}{}}{\displaystyle 𝑑\omega e^{i(k\varphi \omega t)}\frac{}{p}\left[\stackrel{~}{F}_{k1}(p,\omega )\stackrel{~}{F}_{k+1}(p,\omega )\right]}`$ with $$\stackrel{~}{F}_k(p,\omega )F_k(p)𝑑p^{}\stackrel{~}{f}_1(p^{},\omega \frac{kp}{M}),$$ (60) which leads to the equation for the Fourier coefficients: $$\left(\omega \frac{kp}{M}\right)\stackrel{~}{f}_k(p,\omega )=\frac{J}{2}\frac{}{p}\left[\stackrel{~}{F}_{k1}(p,\omega )\stackrel{~}{F}_{k+1}(p,\omega )\right].$$ (61) Since we are dealing with the perturbation of the non-stationary state with a continuous spectrum, the frequency of the perturbation should satisfy $`\omega kp/M0`$; otherwise, there would be no perturbation at all. This allows us to divide Eq. (61) by $`\omega kp/M`$ and to integrate over $`p`$. For $`k=1`$, we have $`{\displaystyle 𝑑p^{}\stackrel{~}{f}_1(p^{},\omega )}`$ $`=`$ $`{\displaystyle \frac{J}{2}}{\displaystyle 𝑑p\left(\omega +\frac{p}{M}\right)^1}`$ (62) $`\times `$ $`{\displaystyle \frac{}{p}}\left[\stackrel{~}{F}_2(p,\omega )\stackrel{~}{F}_0(p,\omega )\right],`$ while for $`k1`$, $`\stackrel{~}{f}_k(p,\omega )`$ is determined by $`\stackrel{~}{f}_1(p^{},\omega \pm kp/M)`$ through Eqs. (60) and (61). It is thus enough to have non-vanishing $`\stackrel{~}{f}_1(p^{},\omega )`$. Now suppose that $`\omega =\omega _0`$ is a solution of Eq. (62), i.e., $`𝑑p^{}\stackrel{~}{f}_1(p^{},\omega )0`$ for $`\omega =\omega _0`$. If we write $`𝑑p^{}\stackrel{~}{f}_1(p^{},\omega )=K\delta (\omega \omega _0)`$, then $`𝑑p^{}\stackrel{~}{f}_1(p^{},\omega +\frac{2p}{M})=K\delta (\omega +\frac{2p}{M}\omega _0)`$. Integration over $`\omega `$ gives $`1+{\displaystyle \frac{JM}{2}}{\displaystyle 𝑑p\frac{F_0^{}(p)}{p+M\omega _0}}`$ $`=`$ $`{\displaystyle \frac{JM}{2}}{\displaystyle 𝑑p\frac{F_2(p)}{(pM\omega _0)^2}}`$ (63) $`=`$ $`{\displaystyle \frac{JM}{2}}{\displaystyle 𝑑p\frac{F_2^{}(p)}{pM\omega _0}},`$ where the last line is obtained by integration by parts. Hence the frequency of a self-sustained oscillation and accordingly the stability is, similarly to the stationary case \[Eq. (43)\], determined by $$1+\frac{JM}{2}𝑑p\left[\frac{F_0^{}(p)}{p+M\omega _0}\frac{F_2^{}(p)}{pM\omega _0}\right]=0.$$ (64) The stability condition is thus entirely the same as that of the stationary case except that we now have two momentum distributions: From Eqs. (85) and (89) with $`M\omega _0=\stackrel{~}{\omega }_r+\stackrel{~}{\omega }_i`$, we have $$\{\begin{array}{cc}\frac{2}{JM}+_{\mathrm{}}^{\mathrm{}}𝑑p\left[\frac{(p+\stackrel{~}{\omega }_r)F_0^{}(p)}{(p+\stackrel{~}{\omega }_r)^2+\stackrel{~}{\omega }_i^2}\frac{(p\stackrel{~}{\omega }_r)F_2^{}(p)}{(p\stackrel{~}{\omega }_r)^2+\stackrel{~}{\omega }_i^2}\right]=0\hfill & \\ _{\mathrm{}}^{\mathrm{}}𝑑p\frac{F_0^{}\left(p\right)}{\left(p+\stackrel{~}{\omega }_r\right)^2+\stackrel{~}{\omega }_i^2}\frac{F_2^{}\left(p\right)}{\left(p\stackrel{~}{\omega }_r\right)^2+\stackrel{~}{\omega }_i^2}=0,\hfill & \end{array}$$ (65) for $`\omega _i>0`$, for which the system is unstable as the perturbation grows in time. In the opposite case ($`\omega _i<0`$), the perturbation dies out to make the system stable. The condition for this is given by $$\{\begin{array}{cc}\frac{2}{JM}+_{\mathrm{}}^{\mathrm{}}𝑑p\left[\frac{(p+\stackrel{~}{\omega }_r)F_0^{}(p)}{(p+\stackrel{~}{\omega }_r)^2+\stackrel{~}{\omega }_i^2}\frac{(p\stackrel{~}{\omega }_r)F_2^{}(p)}{(p\stackrel{~}{\omega }_r)^2+\stackrel{~}{\omega }_i^2}\right]\hfill & \\ +2\pi \left[\text{Im}F_0^{}(\stackrel{~}{\omega })\text{Im}F_2^{}(\stackrel{~}{\omega })\right]=0\hfill & \\ \stackrel{~}{\omega }_i_{\mathrm{}}^{\mathrm{}}𝑑p\left[\frac{F_0^{}(p)}{(p+\stackrel{~}{\omega }_r)^2+\stackrel{~}{\omega }_i^2}\frac{F_2^{}(p)}{(p\stackrel{~}{\omega }_r)^2+\stackrel{~}{\omega }_i^2}\right]\hfill & \\ +2\pi \left[\text{Re}F_0^{}(\stackrel{~}{\omega })\text{Re}F_2^{}(\stackrel{~}{\omega })\right]=0.\hfill & \end{array}$$ (66) Finally, in the neutral case ($`\omega _i=0`$), the condition simply reads $$\{\begin{array}{cc}\frac{2}{JM}+𝒫_{\mathrm{}}^{\mathrm{}}𝑑p\left[\frac{F_0^{}(p)}{p+\stackrel{~}{\omega }_r}\frac{F_2^{}(p)}{p\stackrel{~}{\omega }_r}\right]=0\hfill & \\ F_0^{}(\stackrel{~}{\omega }_r)F_2^{}(\stackrel{~}{\omega }_r)=0,\hfill & \end{array}$$ (67) where $`𝒫`$ stands for the principal part. Our next task is to determine stability for specific distributions of $`F_0(p)`$ and $`F_2(p)`$. Since most dynamical calculations, for both ferromagnetic and antiferromagnetic system, have used the so-called water-bag distribution, we also consider the momenta to be distributed uniformly in the range $`[\alpha ,\alpha ]`$: $$F_0(p)=\pm F_2(p)=\frac{1}{2\alpha }$$ (68) Substitution of $`F_0(p)=\pm F_2(p)=(2\alpha )^1[\delta (p+\alpha )\delta (p\alpha )]`$ into Eqs. (65)-(67), depending on the sign of $`\omega _i`$, determines the frequency $`\omega _0`$. We first consider the case $`F_0(p)=F_2(p)`$. From the second equations of Eqs. (65) (for $`\omega _i>0`$) and (66) (for $`\omega _i<0`$), we find $`\omega _r=0`$, while there is no solution to satisfy the first ones. When $`\omega _i=0`$, again there is no solution to satisfy the first equation of Eq. (67). This indicates that there is no self-sustained oscillation in the system. Note, however, that the system is neutrally stable as it has a continuous spectrum ($`\omega =kp/M`$). Next, when $`F_0(p)=F_2(p)`$, one finds $`\omega _i`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{J}{M}}\left({\displaystyle \frac{\alpha }{M}}\right)^2},\omega _r=0\text{for}\alpha <\alpha _R`$ (69) $`\omega _r`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{J}{M}}+\left({\displaystyle \frac{\alpha }{M}}\right)^2},\omega _i=0\text{for}\alpha >\alpha _R`$ (70) with $`\alpha _R\sqrt{JM}`$. In the microcanonical ensemble one may relate the average kinetic energy with the temperature: $`T/2=p^2/2M=\alpha ^2/6M`$, from which one has $`T_R=\alpha _R^2/3M=J/3`$ . Accordingly, it is concluded in this case that the rotating solution is neutrally stable for $`T>T_R`$ and becomes unstable below $`T_R`$. Note here that $`T_R`$ is lower than the equilibrium critical temperature $`T_c=J/2`$. For the antiferromagnetic system ($`J<0`$), we replace $`J=\|J\|`$ in Eq. (69) to obtain, for $`\omega _i=0`$, $$\omega _r=\pm \sqrt{\frac{\|J\|}{M}+\left(\frac{\alpha }{M}\right)^2},$$ (71) while there is no solution for $`\omega _r=0`$. We therefore conclude that the antiferromagnetic system is neutrally stable for all $`\alpha `$, i.e., at all temperatures. In Sec. 3 we have shown that the bi-cluster state is allowed by the rotating distribution. The result obtained here that this non-stationary solution is neutrally stable at all temperatures thus suggests an alternative explanation as to the origin of the spontaneously formed bi-cluster state in numerical simulations, which retains its form for quite a long time . This keeps parallel with the emergence of quasi-stationarity in the ferromagnetic system, associated with the neutral stability . ## 6 Conclusion In this paper, we have presented a detailed analysis of the system of globally coupled rotors. Starting from a set of Langevin equations and their corresponding FPE, which includes the microcanonical ensemble approach as a limiting case, we have found a class of solutions and studied their stability. The standard canonical distribution constitutes a simultaneous solution of the canonical and the microcanonical ensembles, and thus describes the same equilibrium behavior in both ensembles, leading to the coherent phase (characterized by a nonzero mono-cluster order parameter, i.e., $`\mathrm{\Delta }^{(1)}0`$) below the critical temperature $`T_c`$ in the ferromagnetic system. The stability of the coherent phase is governed by an infinite-order difference equation, the behavior of which may be understood by considering successively higher-order terms (i.e., Fourier components in the perturbation). It has been found that the coherent phase is stable above some temperature $`T_0`$, which is finite if one includes only a few lowest Fourier components. As more components are considered, $`T_0`$ tends to decrease toward zero; this leads us to surmise that the infinite number of Fourier components would stabilize the coherent phase down to zero temperature. Namely, it is expected that the stability equation, if treated exactly, leads to the stability of the coherent phase at all temperatures below $`T_c`$. We find the more interesting possibility for the non-stationary (rotating) solution with regard to dynamical order. Dynamical order, manifested by multi-cluster motion, is allowed for both ferromagnetic and antiferromagnetic interactions. Unlike a ferromagnetic system, in which dynamical order ceases to exist below the temperature $`T_R`$, dynamical order is observed to be neutrally stable down to zero temperature in the antiferromagnetic system. This suggests an alternative explanation as to the origin of the spontaneous formation of the bi-cluster phase in the system of antiferromagnetically coupled rotors. This is in parallel with the explanation that the quasi-stationarity observed in ferromagnetically coupled rotors is related to the neutral stability of the stationary solution in the incoherent phase below the equilibrium critical temperature . To conclude, we have introduced a unified approach for both the canonical ensemble and the microcanonical ensemble, based on the Fokker-Planck equation. Depending on the ensemble, the Fokker-Planck equation admits a few solutions which have implications on some remarkable features (quasi-stationarity in ferromagnetic systems and bi-cluster motion in antiferromagnetic systems) observed in numerical experiments. We provide natural explanations for the origin of those seemingly unrelated features within the same context. Finally, we point out that our approach is based on an effectively one-particle dynamics, exact for an infinite number of particles and does not reflect the instabilities that may be caused by the finiteness of the number of particles. JC thanks the Korea Institute for Advanced Study for hospitality during his stay, where this work was completed. This work was supported in part by the Korea Science and Engineering Foundation through National Core Research Center for Systems Bio-Dynamics and by the Ministry of Education through the BK21 Program. * ## Appendix A Properties of $`\chi _k(\omega )`$ In this appendix we describe some properties of the response function $`\chi _k(\omega )`$ for the Maxwell distribution $`f_M(p)`$: $$\chi _k(\omega )=\frac{J}{2}𝑑p\frac{f_M^{}(p)}{\omega +kp/M}.$$ (72) First, for $`k=0`$, we have $$\chi _0(\omega )=\frac{J}{2}𝑑p\frac{f_M^{}(p)}{\omega }=0$$ (73) since $`f_M^{}(p)`$ is an odd function. We next write $`kk`$ and change the integration variable $`p`$ to $`p`$ in Eq. (72) to get $`\chi _k(\omega )`$ $`=`$ $`{\displaystyle \frac{J}{2}}{\displaystyle 𝑑p\frac{f_M^{}(p)}{\omega kp/M}}`$ (74) $`=`$ $`{\displaystyle \frac{J}{2}}{\displaystyle d(p)\frac{f_M^{}(p)}{\omega +kp/M}}`$ $`=`$ $`{\displaystyle \frac{J}{2}}{\displaystyle 𝑑p\frac{f_M^{}(p)}{\omega +kp/M}}=\chi _k(\omega ),`$ again noting that $`f_M^{}(p)`$ is an odd function. Similarly, it is straightforward to show that $$\chi _k(\omega )=\chi _k(\omega )=\chi _k(\omega ).$$ (75) Further, Eq. (72) can also be written as $`\chi _k(\omega )`$ $`=`$ $`{\displaystyle \frac{JM}{2k}}{\displaystyle 𝑑p\frac{f_M^{}(p)}{p+M\omega /k}}`$ (76) $``$ $`{\displaystyle \frac{1}{k}}\chi (\omega /k),`$ where $$\chi (\omega )\frac{JM}{2}𝑑p\frac{f_M^{}(p)}{p+M\omega }=\chi (\omega )$$ (77) is the response function already defined in Sec. 4. Although we consider here the Maxwell distribution, the properties given above hold for any momentum distribution $`f_0(p)`$, only if it is an even function of $`p`$. We now proceed to evaluate this function, paying attention to the simple pole at $`p=M\omega `$ on the complex $`p`$-plane. Setting $`M\omega \stackrel{~}{\omega }`$ and making analytic continuation $`\omega =\omega _r+i\omega _i`$, we obtain $`\chi (\omega )`$ in the form $`\stackrel{~}{\chi }(\omega ){\displaystyle \frac{2}{JM}}\chi (\omega )=\{\begin{array}{cc}_{\mathrm{}}^{\mathrm{}}𝑑p\frac{f_M^{}\left(p\right)}{p+\stackrel{~}{\omega }}\hfill & \text{for}\omega _i>0\hfill \\ 𝒫_{\mathrm{}}^{\mathrm{}}𝑑p\frac{f_M^{}\left(p\right)}{p+\stackrel{~}{\omega }}i\pi f_M^{}(\stackrel{~}{\omega })\hfill & \text{for}\omega _i=0\hfill \\ _{\mathrm{}}^{\mathrm{}}𝑑p\frac{f_M^{}\left(p\right)}{p+\stackrel{~}{\omega }}2i\pi f_M^{}(\stackrel{~}{\omega })\hfill & \text{for}\omega _i<0.\hfill \end{array}`$ (81) With the tilde sign omitted for convenience, the real part reads $`\mathrm{Re}\stackrel{~}{\chi }(\omega )`$ $`=`$ $`\{\begin{array}{cc}_{\mathrm{}}^{\mathrm{}}𝑑p\frac{\left(p+\omega _r\right)f_M^{}\left(p\right)}{\left(p+\omega _r\right)^2+\omega _i^2}\hfill & \text{for}\omega _i>0\hfill \\ P_{\mathrm{}}^{\mathrm{}}𝑑p\frac{f_M^{}\left(p\right)}{p+\omega _r}+\pi \mathrm{Im}f_M^{}(\omega )\hfill & \text{for}\omega _i=0\hfill \\ _{\mathrm{}}^{\mathrm{}}𝑑p\frac{\left(p+\omega _r\right)f_M^{}\left(p\right)}{\left(p+\omega _r\right)^2+\omega _i^2}+2\pi \mathrm{Im}f_M^{}(\omega )\hfill & \text{for}\omega _i<0\hfill \end{array}`$ (85) while the imaginary part is given by $`\mathrm{Im}\stackrel{~}{\chi }(\omega )`$ $`=`$ $`\{\begin{array}{cc}\omega _i_{\mathrm{}}^{\mathrm{}}𝑑p\frac{f_M^{}\left(p\right)}{\left(p+\omega _r\right)^2+\omega _i^2}\hfill & \text{for}\omega _i>0\hfill \\ \text{Re}f_M^{}(\omega )\hfill & \text{for}\omega _i=0\hfill \\ \omega _i_{\mathrm{}}^{\mathrm{}}𝑑p\frac{f_M^{}\left(p\right)}{\left(p+\omega _r\right)^2+\omega _i^2}+2\pi \mathrm{Re}f_M^{}(\omega )\hfill & \text{for}\omega _i<0.\hfill \end{array}`$ (89) We next write $`f_M^{}(\omega _r+i\omega _i)`$ $`=`$ $`{\displaystyle \frac{\omega _r+i\omega _i}{\sqrt{2\pi M^3T^3}}}e^{(\omega _r^2\omega _i^2)/2MT}e^{i\omega _r\omega _i/MT}`$ (90) $`=`$ $`{\displaystyle \frac{e^{(\omega _r^2\omega _i^2)/2MT}}{\sqrt{2\pi M^3T^3}}}[\omega _r\mathrm{cos}{\displaystyle \frac{\omega _r\omega _i}{MT}}+\omega _i\mathrm{sin}{\displaystyle \frac{\omega _r\omega _i}{MT}}`$ $`+i(\omega _i\mathrm{cos}{\displaystyle \frac{\omega _r\omega _i}{MT}}\omega _r\mathrm{sin}{\displaystyle \frac{\omega _r\omega _i}{MT}})]`$ $``$ $`\mathrm{Re}f_M^{}(\omega _r+i\omega _i)+i\mathrm{Im}f_M^{}(\omega _r+i\omega _i),`$ from which it is obvious that $`\mathrm{Re}f_M^{}(\omega _r+i\omega _i)=0`$ for $`\omega _r=0`$ and $`\mathrm{Im}f_M^{}(i\omega _i)=(2\pi M^3T^3)^{1/2}\omega _ie^{\omega _i^2/2MT}`$. We thus conclude that $`\omega _r=0`$ is a solution of $`\mathrm{Im}\stackrel{~}{\chi }(\omega )=0`$, since $`f_M^{}(p)`$ is an odd function of $`p`$, which makes the integrals vanish in Eq. (89). We now evaluate the integral of $`\mathrm{Re}\stackrel{~}{\chi }(\omega )`$. For $`\omega _r>0`$, the first equation in Eq. (85) becomes $`\text{Re}\stackrel{~}{\chi }(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{T}}[1\sqrt{\pi }ye^{y^2}\text{erfc}(y)]`$ (91) $``$ $`{\displaystyle \frac{1}{T}}f(y)`$ with the scaled variable $`y\omega _i\sqrt{M/2T}`$, where $$\text{erfc}(y)=\frac{2}{\sqrt{\pi }}_y^{\mathrm{}}e^{t^2}𝑑t$$ (92) is the complimentary error function. For $`\omega _i=\|\omega _i\|<0`$, it is straightforward to show that the last equation in Eq. (85) becomes $`\mathrm{Re}\stackrel{~}{\chi }(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{T}}[1+\sqrt{\pi }\|y\|e^{y^2}(2\mathrm{erfc}(\|y\|))]`$ (93) $``$ $`{\displaystyle \frac{1}{T}}g(y).`$ For $`\omega _i=0`$, we have $`\mathrm{Im}f_M^{}(\omega _r+i\omega _i)=0`$ and the second equation in Eq. (85) simply reduces to $`\mathrm{Re}\stackrel{~}{\chi }(\omega )=1/T`$. Note that $`f(y)`$ is a monotonically decreasing function of $`y`$, varying from unity to zero as $`y`$ grows from zero to arbitrarily large values. On the other hand, $`g(y)`$ increases monotonically with $`y`$, from unity to arbitrarily large values. ## References
warning/0506/astro-ph0506657.html	ar5iv	text	# X-Ray Observations of the edge-on star-forming galaxy NGC 891 and its supernova SN1986J ## 1 Introduction The environment of a spiral galaxy can influence its evolution in many ways. Close interactions with nearby galaxies have been shown to trigger inspire star formation, and thus alter the interstellar medium (ISM) of the galaxy. Properties of the ISM could also be affected by direct infall of material from the intergalactic medium (IGM), or, if the galaxy is in a group or cluster, by the ram pressure stripping by the IGM. However, in certain, outflows into the IGM could be a significant mode of the interaction of a galaxy with its surroundings. Bursts of star formation in the galaxy can, if sufficiently strong, lead to the outflow of matter, which will be predominantly visible at H$`\alpha `$ and X-ray wavelengths (Shopbell & Bland-Hawthorn, 1998; Lehnert, Heckman & Weaver, 1999). These outflows can be strong enough to escape the visible extent of the galaxy, but still be retained by the halo of the galaxy. In some of these cases, the burst of star formation would be strong enough for the outflow to be able to could escape the galaxy and its halo, and appear as a superwind. A detailed study of the region of interface between the optical extent of the galaxy and the surrounding medium, is clearly important. Edge-on spiral galaxies, such as the galaxy we study here, NGC 891, are oriented in such a way that these are particularly useful in the study of this interface region (Strickland et al., 2004). Because of the energies associated with supernovae and the wind of massive stars, observations at X-ray energies are of particular interest in studying the hot haloes or outflows from galaxies. Of the two currently operating major X-ray missions, Chandra has very high spatial resolution and moderate collecting area, whereas XMM-Newton has reasonable spatial resolution, but a larger collecting area, and thus should be able to see lower X-ray surface brightness features. In this paper we will present results from XMM-Newton and Chandra observations of NGC 891, an edge-on spiral galaxy. NGC 891 is very similar in many respects to our own galaxy. The high inclination makes NGC 891 ideal for studying the diffuse X-ray emission from the outflow extending from the plane of the galaxy (particularly in the NW direction), which has previously been studied with the ROSAT (Bregman & Pildis, 1992), ASCA (Bregman & Pildis, 1994; Read, Ponman & Strickland, 1997) and Chandra (Strickland et al., 2004) observatories. High-resolution X-ray studies reveal diffuse emission up to about $`10^{37}`$ erg s<sup>-1</sup>, but also reveal a rich point source associated with stellar sources, such as X-ray binaries and supernovae. The X-ray point source population is closely related with the history of star formation in the galaxy (e.g., Kilgard et al., 2002). The slope of the X-ray luminosity function (XLF) of point sources is believed to be a good indicator of this, with flatter XLF slopes indicating more recent star-formation (due to a large fraction of high luminosity massive X-ray binaries). The point source populations of several nearby galaxies have been analysed (Hartwell et al., 2004; Kilgard et al., 2002; Colbert et al., 2004). Of them, Hartwell et al. (2004) selected a small group of well-studied spirals, and found that the XLF slope is generally steeper for normal spiral galaxies than for the starburst galaxies, and correlated with other observable parameters. In this work, we compare these X-ray, optical and IR properties of NGC 891 with a similar sample of normal and starburst galaxies. As part of the point source population, the presence of an active galactic nuclei (AGN) is also the subject of much debate in NGC 891 and other galaxies, and how dominant or irrelevant the AGN is to the energetics of the central regions of the galaxy. Strickland et al. (2004) show that there is the possibility of a weak hard X-ray source (2.0-8.0 keV) using the Chandra data. The discovery of a nuclear AGN in NGC 891 would help us to categorise it as a starburst galaxy. A rare form X-ray point sources in galaxies are X-ray luminous young supernovae. There have been few X-ray detections of young supernovae, and the ones that have been detected are associated with relatively nearby Type II SN events (Houck et al. (1998) and references within). One of these rare objects, detected at X-ray wavelengths, is SN1986J, located in NGC 891 and detected as a bright source in our X-ray observations. It is also very radio bright. Houck et al. (1998) studied the evolution of SN1986J using both ASCA and ROSAT data. SN1986J is still visible at X-ray energies and we present a new analysis of the evolution of this source using both Chandra and XMM-Newton data. The paper is organised as follows: In §2 we discuss the XMM-Newton and Chandra data used here, and present the reduction procedure and results of the analysis of point sources (Chandra) and the diffuse emission (XMM) in §3. In §4, a comparison of X-ray and near and far-IR properties is made with other nearby spiral galaxies. Observations of the supernova SN1986J are presented in §5, and compared with observations at previous epochs. General conclusions are presented in the final section. ## 2 OBSERVATIONS NGC 891 is similar to the Milky Way in optical luminosity ($`m_B^{0,i}=9.4`$ from RC3, implying $`M_B=20.4`$), Hubble type (Sb) and rotational velocity (225 km s<sup>-1</sup>, van der Kruit, 1984; Rupen, 1991). It is a nearby edge-on spiral, and thus has been the subject of many detailed studies of interstellar dust and gas, from observations of the radio continuum (e.g., Dahlem, Dettmar, & Hummel, 1994), H I (e.g., Swaters, Sancisi, & van der Hulst, 1997), carbon monoxide (e.g., Sofue & Nakai, 1993) and molecular hydrogen (e.g., Valentijn & van der Werf, 1999). However, there is considerably more star formation in NGC 891 than in the Milky Way, presumably due to the presence of about 2.5 times as much molecular gas (Scoville et al., 1993). As a result, NGC 891 is found to be twice as luminous at IR wavelengths (Wainscoat, de Jong & Wesselius, 1987), and there is evidence of enhanced dust extinction for optical light, even outside the plane (e.g., Howk & Savage, 1997). The extended gaseous halo of the galaxy in evident from $`H\alpha `$ observations using HST/WFPC2, with filaments reaching up to 2.2 kpc above the galactic plane (Rossa et al., 2004). Other wavelengths have yielded other interesting details. Polarised radio emission has been detected from NGC 891 (Sukumar and Allen, 1991), illustrating the nature of the interstellar magnetic field. A luminous radio halo has been well established with a scale height of 2.7 kpc (Allen, Baldwin & Sancisi, 1978). Throughout this analysis, we assume a distance of $`9.08\pm 0.45`$ Mpc for NGC 891, which is a weighted mean of the values obtained by Ciardullo, Jacoby & Harris (1991) and Tonry et al. (2001)). Other selected parameters are shown in Table 1. Where relevant, we use $`H_0`$= 72 km s<sup>-1</sup> Mpc<sup>-1</sup>. ### 2.1 XMM-Newton observations NGC 891 was observed with the XMM-Newton telescope on August 22 2002, as part of a GTO observation. The galaxy was at the centre of the field of view, with the entire emission also contained within the field of view. The XMM-Newton data was observed with three different cameras, MOS1, MOS2 and PN. Each of the data sets were analysed separately. After running the standard SAS tasks: cifbuild and odfingest; the chains to process the odfs were run (emchain, epchain) simultaneously running the badpixfind command. These have proved to be considerably better at detecting bad pixels than the tasks emproc and epproc. Having extracted the lightcurve in the 10–15 keV energy range, the observation was found t contain a considerable amount of flaring which we used to identify the good time intervals. The periods of high background were identified by eye at 55/40.0 counts for the MOS/PN cameras respectively. The resulting good times are 14.6 ks (originally 18.0 ks) for MOS1/MOS2; and 9.6 ks (originally 15.0 ks) for PN. Events were filtered based on the the pattern parameter, which indicates the geometry of the detection of each event, i.e. the number of adjacent pixels that detect each photon. The events that were retained in the filtering were the ones that are well-calibrated (singles, doubles and quadruples for MOS & singles and doubles for PN). Due to the extent of the diffuse emission present in this galaxy, we used blank sky background files for subtraction of our data. We followed the technique set out by Read & Ponman (2003) for producing images and spectra that are correctly background subtracted. Given the high density of point sources, and the high inclination, it seemed that manually setting a region with exclusion regions based around the point sources (extracted from Chandra data– see below) was more appropriate. The radius of the exclusion region was set by the SAS task calview, using the encircled energy (PSF) function, at $`15^{\prime \prime }`$ (corresponding to a PSF fraction of 75%). When the sources were removed, we used the CIAO task dmfilth to interpolate over the holes. The width of the background annuli were initially taken to be 1.5 times the radius of the source region, but many were adjusted to ensure that the background annuli did not overlap with any source regions. In some cases it was impossible to extract adequate background annuli due to the compact nature of the sources on each other. In that case, the sources were merged and treated as a single source removal region. When analysing the spectral data, the data from the three cameras were fitted simultaneously together to ensure the best fit. However, for the radial profile analysis, the data was mosaicked together using the task emosaic. ### 2.2 Chandra observations To better characterise the point sources, we used the archived Chandra/ACIS-S observation (51 ks, PI: Bregman), which was reduced in a similar way, using the CIAO software with online threads. Once the Chandra data preparation was completed, the CIAO tool wavdetect was run on the Chandra data. We narrowed down the returned list of sources by excluding all the sources outside the $`D_{25}`$ ellipse (see Table 1), and considering only the sources inside. We used the wavelet scales 1, 2 and 4 to extract the point sources, as these provided the best identification of all the point sources present. These sources are shown on a DSS overlay in Fig. 1). We used the results to generate an X-ray flux of each source to generate a luminosity function (Fig. 2). ## 3 ANALYSIS OF THE X-RAY OBSERVATIONS Here we analyse the X-ray observations described in the previous section to characterise the population of point sources in NGC 891, using the archived Chandra data set, utilising the superior resolution of Chandra/ACIS setup. We use this observation to subtract the point sources from the XMM-Newton observation, to investigate the nature of the diffuse emission. ### 3.1 The Point Source Luminosity Function from Chandra We detected a total of 26 point sources within the $`D_{25}`$ ellipse, down to a flux limit of $`10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, which at the adopted distance of NGC 891 amounts to a luminosity of $`F_{min}=10^{37}`$ erg s<sup>-1</sup>. The regions were all output as ellipses containing 99.7% (3$`\sigma `$) of the source counts. As a considerable number of the sources contained few counts, an accurate spectral fit was not possible, so we modelled each of the point sources individually in XSPEC using an absorbed power-law model with a slope of $`\mathrm{\Gamma }=1.8`$. The point sources with more than 100 counts were individually fitted. Flux values were established for each of the point sources in an energy range of 0.3–8.0 keV, and these were combined to make a luminosity function (Fig. 2). We discuss the point source associated with SN1986J in greater detail in the following section. As we will see below, the slope of the luminosity function can be used as an indicator of recent star formation activity in a galaxy. This so-called $`\mathrm{log}N`$-$`\mathrm{log}S`$ relation is usually expressed in the cumulative form $$\mathrm{log}N(>S)=\alpha \mathrm{log}S+\kappa $$ (1) where a straight line is fitted to the cumulative histogram. Here we adopt a more robust method of measuring the slope of the luminosity function with the potential of better measuring the uncertainty of the slope. This is a generalised version of the much-used method adapted from Crawford, Jauncey, & Murdoch (1970) and Murdoch, Crawford, & Jauncey (1973), the crucial difference being that in the former, errors in flux measures are not taken into account and in the latter, all errors are considered equal. Here, each of our $`n`$ point sources has a measured luminosity $`F_i`$, given the adopted distance to the galaxy, and an independently estimated error $`\sigma _i`$. We represent the probability distribution of the luminosity $`S`$ of point sources in the galaxy in its differential form $$P(S)dS=AS^\beta dS.$$ (2) On comparison with (1), $`\alpha =\beta 1`$. Our exercise thus consists of maximising the log likelihood function $$=\underset{i=1}{\overset{n}{}}\mathrm{ln}P(F_i,\sigma _i),$$ (3) where the distribution of our measured values of flux and standard deviation $`(F_i,\sigma _i)`$ is given by $$P(F_i,\sigma _i)=\frac{_0^{\mathrm{}}P(F_i,\sigma _i\|S)P(S)𝑑S}{_{F_{\mathrm{min}}}^{\mathrm{}}_0^{\mathrm{}}P(F_i,\sigma _i\|S)P(S)𝑑S𝑑F}.$$ (4) Assuming the errors of measuring flux and luminosity are distributed as a Gaussian, the integrand above is given by $$P(F_i,\sigma _i\|S)P(S)dS=\frac{A}{\sigma _i\sqrt{2\pi }}S^\beta e^{(F_iS)^2/2\sigma _i^2}dS.$$ (5) We numerically find the value of $`\beta `$ (thus $`\alpha `$) for which $``$ in (3) is maximum. For large $`N`$, the probability distribution for $`\beta `$ is asymptotically Gaussian, and the $`1\sigma `$ error in $`\beta `$ corresponds to $`\mathrm{\Delta }=0.5`$. For our limiting 0.3-8 keV luminosity, $`F_{min}=10^{37}`$ erg s<sup>-1</sup>, we find the slope of the cumulative LF given by (1) to be $`\alpha =0.77_{0.10}^{+0.13}`$, as shown in Fig. 2. Since NGC 891 is almost exactly an edge-on galaxy, the effect of differential local extinction might be important within the point source population, which might lead us to progressively miss the fainter objects on the far side of the galaxy. Since the effect of such extinction is energy-dependent (being larger at lower energies, particularly below 2 keV), we tested whether it is substantial by re-calculating the luminosity function for the energy range 2–8 keV (Fig. 2b). The slope was found to be the same at $`\alpha =0.77_{0.12}^{+0.16}`$, leading us to conclude that the effect of differential extinction is not very serious for this target. ### 3.2 The diffuse X-ray emission from XMM-Newton We extracted the spectrum of the diffuse emission in the energy range 0.3–6 keV, excluding a circle of 15 arcsec around each detected point source, from the 3 XMM-Newton camera data and analysed them with the heasoft package XSPEC. We performed simultaneous fitting of all three data sets in order to fit the best model. We also tried to add the Chandra diffuse emission to the fits, but due to the limited collecting area of the Chandra telescope, we found the diffuse emission spectra generated did not improve the fit. We initially fitted an absorbed single temperature plasma model with a power-law component to the data. We fixed the Galactic value of column density $`n_H`$ of 6.78$`\times `$10<sup>20</sup> cm<sup>-2</sup> (obtained from Colden). Since the target is an edge-on galaxy, the local column density is expected to be substantial- this was looked for, and fixed at a fairly low value of $`n_H`$ (local) at 2.0$`\times `$10<sup>19</sup> cm<sup>-2</sup>. The power-law model was added to our data to account for point source contamination. Figure 3a shows the best fit to all three data sets, having a $`\chi ^2`$=68.35 with 65 d.o.f., leading to $`\chi _{red}^2`$=1.06. We extract a value of $`kT=0.26\pm 0.01`$ keV. In comparison, Read, Ponman & Strickland (1997) fit a single temperature model to the data with a cooler temperature component ($`0.11\pm 0.03`$ keV). Bregman & Pildis (1994) found that the data was best fitted by a two-temperature model with $`kT`$ at $`0.31`$ and $`10`$ keV. As the result from Bregman & Pildis (1994) implied a two temperature model, a second temperature component was added to the data (Fig. 3b). The two temperatures that were inferred, as a result, were $`kT=0.08\pm 0.01`$ keV and $`0.30\pm 0.03`$ keV. The two temperatures agree within errors to the cooler temperature component derived by Bregman & Pildis (1994) and the single temperature component derived by Read, Ponman & Strickland (1997). However, it should be stated that the spectrum of the data was best constrained by a single temperature fit, and the improvement in the $`\chi _{red}^2`$ by adding a second temperature component is not statistically significant (f-test statistic value $`=1.05`$). A summary of our fit parameters compared with those of Bregman & Pildis (1994) and Read, Ponman & Strickland (1997) are presented in Table 2. Our results do seem to complement the ROSAT data to some extent. Our single and dual temperature model values are consistent with temperatures that both have derived. Bregman & Pildis (1994) do find a high temperature component at 10 keV, which is not determined from our fits. The reason for this may be the poor resolution of ROSAT, which means that it is difficult to remove point sources effectively. This effect is compounded by the high inclination of the system. Also, SN1986J was considerably brighter, at the time of the ROSAT observation more than ten years ago, than it is now, and may well have contaminated the extracted diffuse emission. Finally, the hard X-ray flux of a nearby unresolved bright point source near SN1986J might also account for this discrepancy (Read, Ponman & Strickland, 1997). We also extracted images of the diffuse emission from the XMM-Newton data. These images complement the X-ray observations from ROSAT (Bregman & Pildis, 1994) and Chandra (Read, Ponman & Strickland, 1997). There is a considerable amount of emission from the centre of the galaxy in the NW direction. Figure 1a shows this result for the soft energy band (0.3–2.0 keV). The X-ray contours are overlaid onto the optical Digitised Sky Survey (DSS) image. The image was smoothed with an adaptive gaussian program asmooth such that the $`1\sigma `$ width of the gaussian contained $`40/3`$ counts around each pixel. We also took a minor axis diffuse emission profile to show the extent over which the extended emission protrudes (Fig. 4). The emission protrudes out to just over 6 kpc from the centre of the galaxy in the NW direction, but drops off rapidly in the SE direction outside the plane of the galaxy. ### 3.3 Evidence of central AGN Using the Chandra observation, Strickland et al. (2004) found a weak hard (2.0-8.0 keV) X-ray source near the nucleus of the galaxy, defined as the position of the radio continuum point source (Rupen, 1991). This was suggested to be a weak AGN possibly associated with the radio source. Unfortunately, we were not able to detect this source with our XMM-Newton observation, most likely due to the poorer resolution compared to Chandra. ## 4 Comparison with other local galaxies In order to characterise the global properties of NGC 891 and find whether its properties are consistent with that of a starburst galaxy or not, we compile various X-ray, near-IR and far-IR properties of several nearby spiral galaxies, of which some are normal galaxies and some are actively star-forming and starburst galaxies. Starting from a sample of galaxies selected from Hartwell et al. (2004), a few galaxies from Strickland et al. (2004) and Hartwell et al. (2004). A selected list of parameters are collected in Table 6. #### Far-Infrared luminosities We calculated the FIR lumiosities (erg s<sup>-1</sup>) differently depending on whether the source was from ISO (preferably) or IRAS. We used the equations from Bendo et al. (2002) (fluxes in Jy): $$L_{FIR}=1.89\times 10^{14}(1.36f_{60}+0.958f_{100}+0.439f_{180}),$$ (6) and for IRAS $$L_{FIR}=1.26\times 10^{14}(2.58f_{60}+f_{100}),$$ (7) #### Near-Infrared luminosities To calculate the K-band luminosity $`L_K`$ (erg s<sup>-1</sup>), we derived the following expression using the zero point values from Cohen, Wheaton & Megeath (2003) : $$\mathrm{log}L_K=43.130.4K_{tot}+2\mathrm{log}D,$$ (8) with $`D`$ in Mpc, and $`K_{tot}`$ the total K-band magnitude taken from the 2MASS galaxy atlas (Jarrett et al., 2003). #### X-ray luminosities $`L_X`$ was established from the literature, but different authors used varying energy ranges. Strickland et al. (2004) evaluate their fluxes for an energy range 0.3–2.0 keV (soft X-ray) We have adopted this energy range to keep our values for $`L_X`$. In the case where the data was not in the necessary energy range, we used the Portable Interactive Multi-Mission Simulator (PIMMS) to generate our data. Using the model fit parameters for the diffuse emission in the literature, we estimated a revised flux in the appropriate energy range. #### Deprojected area of galaxy We differ from the calculation made by Hartwell et al. (2004) for the area of the $`D_{25}`$ ellipse. Hartwell et al. (2004) calculated the galaxy area as $`\pi ab\mathrm{cos}\theta `$. For systems that are face on (such as NGC 4214) this calculation is adequate. However, for very high inclination systems, the measure of the semi-minor axis simply represents the width of the disk, and cannot be used in calculating the area. Therefore, for the sake of consistency we choose to evaluate our area as $`\pi a^2`$ for all galaxies in our sample. ### 4.1 Starburst or not? There has been much discussion as to whether NGC 891 is a normal spiral galaxy, or whether it is a starburst (Strickland et al., 2004). We try to categorise NGC 891 by means of comparison with a sample of 16 nearby galaxies. Figure 5 contains four plots, which are described below. The slope of the luminosity function can be used as an indicator of the general properties of the host galaxy. If the star formation rate in the galaxy is constant, then the luminosity function of the sources should be fit by a single unbroken power-law model. If the galaxy is a starburst one, then new high mass X-ray binaries (HMXBs) would be formed, breaking the luminosity function slope. The break in the slope would decrease with time, and would be an indication of the time of previous bursts in the galaxy. As we saw above, in NGC 891, we chose to fit the XLF slope can be fitted by a single power-law. In the case of NGC 1482, the galaxy is further away (22 Mpc) and only two sources could be detected so no XLF slope was determined. Two other galaxies NGC 4244 and NGC 6503 also have low statistics (with 3 and 4 sources respectively), but XLF slopes were calculated by Colbert et al. (2004). We include these in our plots, but they do not appear to be representative of the normal spirals category. Figure 5a plots the XLF against the $`S_{60}/S_{100}`$ flux ratio. Lehnert & Heckman (1996) analyse the properties of superwinds and quantitatively discuss the likelihood of generating a superwind with respect to the infra-red properties of the galaxy. The galaxies with high luminosity ($`L_{IR}10^{44}`$ erg s<sup>-1</sup>), large infra-red excesses ($`L_{IR}/L_{OPT}2`$) and warm far infra-red colours ($`S_{60}/S_{100}0.5`$) are most likely to have galactic superwinds. For the sample selected, the last criterion is satisfied for all the starburst galaxies. Of the starburst galaxies in our sample, superwinds are detected in M82 (Strickland et al., 1997), NGC 253 (Strickland et al., 2000) and NGC 4945 (Moorwood et al., 1996). NGC 5253 (Strickland & Stevens, 1999) and NGC 4449 (Summers et al., 2003) exhibit emission from superbubbles. Although superwinds and superbubbles are both the result of massive star formation processes within the densest regions of the hot galaxies, the winds are able to channel the metals produced straight into the IGM, whereas the superbubbles do not reach the outskirts of the host galaxies, hence retaining their newly processed metals, consequently raising the abundance of the ISM (Tenorio-Tagle, Silich & Mu$`\stackrel{~}{n}`$oz-Tu$`\stackrel{~}{n}\stackrel{`}{o}`$n, 2003). The only starburst galaxies not to have superwinds or superbubbles are NGC 4214, which is calculated likely to exhibit blowout (Hartwell et al., 2004) and the Antennae, which is an interacting system, so no superwind has formed due to disruption from the merging. Therefore, there are two distinct categories of result: The low XLF slope galaxies have a higher ratio of warm IR colours than the high XLF slope galaxies, which are normal spiral galaxies. NGC 891 appears to have the XLF slope of a starburst, but the $`S_{60}/S_{100}`$ corresponding more to normal spirals than starbursts. Figure 5b plots the far infra-red luminosity ($`L_{FIR}`$) against the XLF slope of the data. $`L_{FIR}`$ is an indicator for how much dust there is in a galaxy, and hence the star formation rate (Kennicutt, 2003). The FIR also contains the most important cooling lines in the neutral ISM. Therefore, if the XLF slope is dependent on star formation rate, then, one would expect a correlation between these two parameters. However, this is not observed in our sample. Both the starbursts and normal galaxies span a wide overlapping range of FIR luminosities. The division between spirals and starbursts is quite clear in this plot, breaking either side of 1. Therefore, an XLF-$`L_{FIR}`$ dependence cannot be justified on the basis of these results. Figure 5c plots the XLF slope against the $`S_{60}`$ flux scaled against the area of the $`D_{25}`$ ellipse.The x-axis therefore represents the star formation rate per unit area. There is a definite correlation between the data points on this plot. The higher star formation rate corresponds to starburst galaxies, which we would expect to see. The major exceptions to this trend are NGC 4244 and M83. The XLF slope of NGC 4244 was determined with only three data points (Colbert et al., 2004), and has an estimated error on the data. M83 should be classified as a starburst, as observations indicate that there are massive clusters of stars in the nuclear region (Harris et al., 2001). However, as discussed in Kilgard et al. (2002) the star formation rate is low compared to other starbursts, and the starburst region is confined to a small area of the galaxy. Therefore, most of the point sources analysed were taken from a region outside of the starburst region, so the XLF slope is more representative of the disk population of normal spirals as opposed to the standard starburst population. In Fig. 5d, the warm IR colour ratio S<sub>60</sub>/S<sub>100</sub> is plotted against the ratio of $`L_{FIR}/L_K`$. This ratio tells us the extent of the star formation activity normalised against the stellar types present in the galaxies. The different galaxy types are clearly grouped separately, and all the galaxies seem to follow a linear trend. The normal spiral galaxies have lower star formation rates, which would imply that they are less likely to generate galactic superwinds. The galaxies with higher star formation rates are all starburst galaxies with superwinds (or superbubbles) present. Again, NGC 891 seems to sit between the two categories, making it difficult to categorise as either a spiral or as a starburst. ### 4.2 The X-ray Schmidt law In a classic paper, Schmidt (1959) showed that the rate of Population I star formation is proportional to the density of available gas in a galaxy. This has been re-discovered in various guises over the last few decades. In terms of the X-ray luminosity of the galaxy, which is related to a density of hot gas, one expects it to correlate with the star formation rate, which we can characterise by the total far-IR flux. Ranalli, Comastri & Setti (2003) analyse the effectiveness of the 2–10 keV X-ray luminosity as a star formation rate indicator, showing the linear relation between the X-ray luminosity and the far-IR luminosities, although there are discrepancies between many other published results. Previous studies using ASCA, ROSAT/PSPC or EINSTEIN values would overestimate $`L_X`$, being unable to separate out the emission from X-ray binaries and AGN. Griffiths & Padovani (1990) established $`L_XL_{FIR}^{0.6}`$, although $`S_{60}`$ made up the x-axis, and $`L_X`$ ranged from 0.5–3.0 keV. Also with the EINSTEIN satellite, David, Jones & Forman (1992) obtained an almost linear fit to their data $`L_XL_{FIR}^{0.95\pm 0.06}`$ agreeing quite well with our results within errors. This result was determined in the 0.5–4.0 keV energy band and the sample was a selection of starburst and normal galaxies. Ranalli, Comastri & Setti (2003) analyses data obtained from the ASCA satellite and for an soft X-ray energy range of 0.5–2.0 keV obtains $`L_XL_{FIR}^{0.87\pm 0.08}`$, an identical result to ours for a very similar energy range (we use 0.3–2.0 keV). We plot $`L_X`$ (the diffuse X-ray luminosity) against $`L_{FIR}`$ for our sample galaxies in Fig. 6. We establish that the relation between the data is $`L_XL_{FIR}^{0.87\pm 0.07}`$, without the two data points on the plot which do not appear to lie near the line. If these are included, the fit becomes considerably worse: $`L_XL_{FIR}^{0.71\pm 0.12}`$. These are the two starburst galaxies NGC 253 and NGC 4945, both of which contain a superwind emitting from the centre. This would indicate the presence of an excess of dust in these galaxies containing considerable star formation. NGC 891 lies comfortably near the best-fitting line, away from the normal galaxies since its X-ray and far-IR fluxes are higher, again showing that it does not have as extreme a rate of star formation as a starburst galaxy. ## 5 The resident young X-ray supernova SN1986J Even though thousands of supernovae have been found in the optical, there are instances of only 15 supernovae being detected in the X-ray (Immler & Lewin, 2003; Bregman et al., 2003). Of these, a handful have been monitored over many years, their X-ray emission dominating their radiative output after about a year after explosion (Pooley et al., 2002). Thermal emission from the reverse shock region is expected to be seen as softer X-ray emission within the expanding SN shell as it interacts with the dense stellar wind of the progenitor (Chevalier & Fransson, 1994). One of these Type II supernovae observed in the X-rays is SN1986J in NGC 891, which was discovered using radio observations from the VLA (Rupen et al., 1987). The first X-ray observations of SN1986J were taken with ROSAT (Bregman & Pildis, 1992). Since then, two more ROSAT and two ASCA observations have been undertaken in the mid 1990’s (Houck et al., 1998). The X-rays are thought to be caused by a shock propagating through the former stellar envelope (Bregman & Pildis, 1992); or due to a the supernova envelope colliding with a shocked clumpy wind, which produces the X-rays (Chugai, 1993). Bregman & Pildis (1992) found that the X-rays were emitted primarily in the soft X-ray band (0.1–2.5 keV). An analysis of the ROSAT and ASCA observations was done by Houck et al. (1998), who determined that the X-ray light curve of the data was dropping by $`t^2`$, which was considered to be a rather fast decline. A summary of all observations are shown in Table 3. Note that the X-ray fluxes derived for the ASCA data are an average over all their different models, and so the error shown here is simply the standard deviation of these data sets. ### 5.1 The spectrum of SN1986J The spectrum for the Chandra observation was extracted using a local background and we fitted different temperature models to the data. Firstly, an absorbed single temperature thermal plasma model (mekal) was fitted with a local and Galactic absorbing column (Bregman & Pildis, 1992). The Galactic value for $`\mathrm{log}n_H`$ was fixed at $`20.8`$ ($`n_H`$ in cm<sup>-2</sup>) and the local absorbing column was allowed to vary. The abundance was initially fixed at 0.3 then 1.0 solar, and finally was allowed to vary. The local $`\mathrm{log}n_H`$ was evaluated at 21.6, in good agreement with the value of Bregman & Pildis (1992) (21.7), but the temperature range that we obtained from the model was 3.5-4.1 keV. This temperature component is much larger than the 2.0 keV values obtained from the ROSAT observations, but smaller than the ASCA temperature ranges. A vmekal model was also tried, where the individual element abundances are allowed to vary, as the mekal fit results were leaving several emission lines unidentified. The result from the vmekal fit is plotted in Fig. 7. The temperature was determined to be 4.6 keV, local $`\mathrm{log}n_H`$=21.5, and X-ray flux 0.45$`{}_{0.7}{}^{+0.1}\times 10^{14}`$ erg s<sup>-1</sup>cm<sup>-2</sup>. Many elements were identified in the spectrum of the supernova, most notably, the Fe emission line at 6.4 keV. The abundances are tabulated in Table 4. On inspection of Fig. 7, some of the emission lines are not fitted that well by a vmekal model. Several methods were tried to improve the fit, including a two temperature model (mekal+vmekal) and a different single temperature thermal plasma model (apec). In each of these models, the hard temperature and local $`n_H`$ were similar to the results obtained from the single temperature fitting. For the two temperature fitting, the lower temperature component was 0.11 keV which is likely the component from the diffuse emission of the galaxy, and not the supernova. The results from all the Chandra fits are tabulated in Table 5. None of the fits identified any of the other emission lines. As all the flux values generated from the fits were similar, the value adopted for our supernova flux is the weighted mean of the different models: 9.6$`{}_{0.4}{}^{+0.6}\times 10^{14}`$ erg s<sup>-1</sup>cm<sup>-2</sup>. The energy range chosen was 0.5-2.5 keV to be consistent with the energy ranges used in the ASCA and ROSAT data. For the XMM-Newton data, the flux was much harder to determine. As the flux is dropping every year, the supernova becomes harder to model the later the observations are. We fitted all three cameras simultaneously to try and obtain the best fit to the data. Our data was best fitted with a single temperature model (mekal) with a temperature of 3.6 keV and absorbing column density $`\mathrm{log}n_H=21.53`$. This fit resulted in a flux of $`8.31\pm 1.06\times 10^{14}`$ erg s<sup>-1</sup> cm<sup>-2</sup>. Fitting a dual temperature model did not improve the fit. The temperatures obtained from the fit were 0.21 keV and 3.82 keV for the soft and hard components respectively. This fit did produce a surprisingly large local $`n_H`$ of 22.6. The two-temperature flux was 8.40$`\pm 1.16\times 10^{14}`$ erg s<sup>-1</sup>cm<sup>-2</sup>. The value adopted from the flux is the mean of the four different temperature fits for the XMM-Newton data: 8.5$`\pm 0.15\times 10^{14}`$ erg s<sup>-1</sup>cm<sup>-2</sup>. ### 5.2 The evolution of SN1986J with time SN1986J has been observed with ROSAT, ASCA, Chandra and XMM over the last twelve years. In Fig. 8 we plot the 0.5–2.5 keV flux against time over this period. We recall that from the ROSAT and ASCA data (Houck et al., 1998), the decline of this flux was found to be $`t^2`$, which was considered to be rather steep. Here we attempt to find a mean rate of decline over a much larger time scale. The data from ASCA and ROSAT appear to be inconsistent (Houck et al., 1998) and so we have established a best fit line inclusive and exclusive of the ASCA data. The slopes are calculated as $`2.99\pm 0.45`$ and $`2.89\pm 0.19`$ respectively. In Fig. 8, we also show an alternative set of points for ASCA where we attempt to correct for the poor resolution by subtracting from the the flux of the SN the combined flux of all the other point sources that are expected to be within its point-spread function. However, our results imply $`L_Xt^3`$. We can compare this with the expected behaviour from the work of Chevalier & Fransson (1994), who show that $`L_Xt^1`$ when free-free emission dominates ($`T>4\times 10^7`$K), and $`L_Xt^{0.7}`$ when line emission dominates ($`10^5<T<4\times 10^7`$K). Such a trend in the decline in luminosity is seen in various other Type II SNe, including SN 1993J in M81 and SN1999em in NGC 1637 (Immler & Lewin, 2003). However, our measured temperature of the plasma is around 4 keV, which corresponds to the former range, which indicates that the decline of $`L_X`$ is unusually steep in this case. It is likely that the poor resolution of the ASCA and ROSAT satellites might not have been able to fully eliminate point source contamination, thus overestimating the fluxes of the data. ## 6 Conclusions In this paper, we have presented a GTO XMM-Newton observation of NGC 891, a nearby edge-on spiral galaxy. We find that the diffuse X-ray emission protrudes from the disk in the NW direction out to approximately 6 kpc. The extraplanar hot gas was previously observed in ROSAT observations, but the XMM observation seems to show a sharp cut-off to this gas, showing that this might be the extent to which the gas has reached. We also find that the best fitting thermal plasma model seems to require two temperatures of 0.08 and 0.3 keV respectively, though the fit of a single-temperature plasma of 0.26 keV isn’t much worse. This contradicts previous findings based on ROSAT observations (Bregman & Pildis, 1994; Read & Ponman, 2003). We use an archived Chandra observation to study the point source population within the $`D_25`$ ellipse of the galaxy. We use a robust maximum likelihood method to determine the slope of the cumulative luminosity function $`N(>S)=S^\alpha `$, and find that the slope is rather shallow, $`\alpha =0.77_{0.10}^{+0.13}`$. We have verified that this isn’t predominantly due to extinction, by plotting the XLF for the energy range $`>2`$ keV, without much change on slope. Using a sample of other local galaxies, we have compared the X-ray and infrared properties of NGC 891 with those of nearby ’normal’ and starburst spiral galaxies. We conclude that NGC 891 has more abundant star formation than a normal spiral, but does not have as extreme properties as starburst galaxies like NGC 253 and NGC 4945, and that it is most likely a starburst galaxy in a quiescent state. We examine an X-ray version of the “Schmidt Law”, which correlates the rate of star formation in a galaxy to the mass or density of its available gas. We show that the diffuse X-ray luminosity, an indicator of gas mass, of nearby spirals scales with their far infra-red luminosity, as indicator of dust mass and star formation rate, as $`L_XL_{FIR}^{0.87\pm 0.07}`$, except for extreme starbursts, and NGC 891 does not fall in the latter category. We study the supernova SN1986J in NGC 891 in both XMM-Newton and Chandra observations, nearly twenty years after its explosion. It was studied in the X-ray using ROSAT and ASCA almost a decade ago. The flux and temperature calculated in those studies of low spatial resolution could have significant contamination from their point sources and the diffuse emission. A direct comparison shows that the temperature of the SN remnant at the time of observation was $`5\times 10^7`$ K, which is less than it was a decade ago, but it also reveals that the X-ray luminosity has been declining with time ($`L_Xt^3`$) far more steeply than expected. ## Acknowledgements Our thanks to Trevor Ponman and Andrew Read for sharing their XMM GTO data with us, and their advice. We thank Ben Maughan, Irini Sakelliou and Ed Colbert for useful comments and help in data analysis, and Joel Bregman for a very interesting discussion. We also thank an anonymous referee for useful comments. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation.
warning/0506/gr-qc0506081.html	ar5iv	text	# Quantum Roots in Geometry: I ## 1 Introduction Most of the success of physics in the 20th century has been achieved as a result of the applications of two philosophies.The first is the Quantization Philosophy and the second is the Geometerization Philosophy. The consequences of applying the first is the Quantum Theory, while the consequences of applying the second is the General Theory of Relativity, (GR). The study and understanding of the four known fundamental interactions are not equally successful using, only, one of these two rival philosophies. Electromagnetism, weak and strong interactions are well understood using the quantization philosophy, while gravity is not understood using this philosophy. In the context of geometerization of physics, GR is considered as a good theory for gravity, while there are no such successful geometric theories for the other three interactions. It seems that a third philosophy is needed to unify the physics of the four fundamental interactions. This philosophy may lead to new physics. This would be, undoubtedly, a difficult task. It would be of importance to reach the conclusion that the two rival philosophies are completely exhausted, before trying a third one. This my be a less difficult task. It needs a careful examination of applying the existing philosophies. Examination of the geometric approach to physics shows that this approach is not exhausted yet. Some types of geometry admit some quantum properties. This is what I am going to show in the present work. The following statement summarizes the philosophy of gemeterization of physics: ”To understand nature one should start with geometry and end with physics”. In applying this philosophy, one should look for an appropriate geometry. Einstein, in applying his geometerization philosophy, used three types of geometry. Some of the main properties of these geometries are summarized in the following table. Table I: Comparison between 3-types of geometry \| Geometry \[Ref.\] \| Metric \| Connection \| Building Blocks (#) \| \| --- \| --- \| --- \| --- \| \| Riemannian \| Symmetric \| Symmetric \| Metric tensor (10) \| \| Absolute Parallelism \| Symmetric \| Non-symmetric \| Tetrad vectors (16) \| \| Einstein Non-symmetric \| Non-symmetric \| Non-symmetric \| Metric tensor (16) \| We mean by the term ”Building Block” the geometric object, using which one can construct the whole geometry. In the last column of Table I, we assume that the dimension of space $`n=4.`$ Riemannian geometry has been used by Einstein to construct his successful theory of gravity, GR. It is well known that the number of building blocks in this geometry is just sufficient to describe gravity. For this reason, we are going to consider the other two geometries, in Table I, since the number of building blocks in each is enough to accommodate other interactions, together with gravity. These interactions may have some quantum properties. The term ”Non-Symmetric Geometry” will be used to indicate that the geometry admits non-symmetric connection. In such a geometry, one can define define three types of tensor derivatives (derivatives that preserve tensor properties): $$A_{+\|\nu }^\mu \stackrel{def.}{=}A_{,\nu }^\mu +A^\alpha C_{.\alpha \nu }^\mu ,$$ $`(1.1)`$ $$A_{\|\nu }^\mu \stackrel{def.}{=}A_{,\nu }^\mu +A^\alpha C_{.\nu \alpha }^\mu ,$$ $`(1.2)`$ $$A_{0\|\nu }^\mu \stackrel{def.}{=}A_{,\nu }^\mu +A^\alpha C_{.(\alpha \nu )}^\mu ,$$ $`(1.3)`$ where $`A^\mu `$ is an arbitrary vector and $`C_{.\nu \alpha }^\mu `$ is the non-symmetric connection. Braces ( ) are used for symmetrization and brackets \[ \] will be used for anti-symmetrization. The comma is used for ordinary (not tensor) partial differentiation. Now, what is the starting point for examining non-symmetric geometries to look for any quantum features? It is well known that, quantum properties in microscopic world were discovered when Planck tried to interpret black body radiation, a phenomena which is closely connected to the motion of electrons. On the other hand, in the context of geometerization of physics, motion is described using paths (curves) of an appropriate geometry. So, a good starting point, may be a search for path equations in the geometries under consideration. Bazanski , has established a new approach to derive the equations of geodesic and geodesic deviation simultaneously by carrying out variation on the following Lagrangian: $$L_B=g_{\mu \nu }U^\mu \frac{D\mathrm{\Psi }^\nu }{DS},$$ $`(1.4)`$ where $`U^\mu \stackrel{\mathrm{def}}{=}\frac{dx^\mu }{dS}`$, $`g_{\mu \nu }`$ is the metric tensor, $`\mathrm{\Psi }^\mu `$ is the deviation vector and $`\frac{D}{DS}`$ is the covariant differential operator using Christoffel Symbol. We are going to generalize the Bazanski approach, by replacing the covariant derivative, used in his Lagrangian, by tensor derivatives of the types given by (1.1), (1.2) and (1.3), admitted by the geometry under consideration. The work in this review is organized as follows: Section 2 gives a brief review of the two non-symmetric geometries under consideration, together with the new path equations resulting from each one. Section 3 gives some remarks about the quantum features appearing in these geometric paths. A method for diffusing the quantum properties, in the whole geometry, is given in Section 4. The general quantum path, of the absolute parallelism geometry, is linearized in Section 5. Section 6 gives confirmation and applications of the quantum paths. the work is discussed and concluded in Section 7. ## 2 Geometries with Built-in Quantum Roots ### 2.1 The Absolute Parallelism Geometry Absolute parallelism (AP)space is an n-dimensional manifold each point of which is labelled by n-independent variables $`x^\nu (\nu =1,2,3,\mathrm{},n)`$ and at each point we define n-linearly independent contravariant vectors $`\underset{i}{\overset{\mu }{𝜆}}(i=1,2,3,3\mathrm{},n`$, denotes the vector number and $`\mu =1,2,3\mathrm{},n`$, denotes the coordinate component) subject to the condition, $$\underset{i}{\overset{\stackrel{\mu }{+}}{𝜆}}{}_{.\|\nu }{}^{}=0,$$ $`(2.1)`$ where the stroke and the (+) sign denote absolute differentiation, using a non-symmetric connection to be defined later. Equation (2.1)is the condition for the absolute parallelism. The covariant components of $`\underset{i}{\overset{\mu }{𝜆}}`$ are defined such that $$\underset{i}{\overset{\mu }{𝜆}}\underset{i}{𝜆}{}_{\nu }{}^{}=\delta _\nu ^\mu ,$$ $`(2.2)`$ and $$\underset{i}{\overset{\nu }{𝜆}}\underset{j}{𝜆}{}_{\nu }{}^{}=\delta _{ij}.$$ $`(2.3).`$ Using these vectors, the following second order symmetric tensors are defined: $$g^{\mu \nu }\stackrel{\mathrm{def}}{=}\underset{i}{\overset{\mu }{𝜆}}\underset{i}{\overset{\nu }{𝜆}},$$ $`(2.4)`$ $$g_{\mu \nu }\stackrel{\mathrm{def}}{=}\underset{i}{𝜆}{}_{\mu }{}^{}\underset{i}{𝜆}{}_{\nu }{}^{},$$ $`(2.5)`$ consequently, $$g^{\mu \alpha }g_{\nu \alpha }=\delta _\nu ^\mu .$$ $`(2.6)`$ These second order tensors can serve as the metric tensor and its conjugate of Riemannian space, associated with the AP-space, when needed. This type of geometry admits, at least, four affine connections. The first is a non-symmetric connection given as a direct solution of the AP-condition(2.1), i.e. $$\mathrm{\Gamma }_{.\mu \nu }^\alpha =\underset{i}{\overset{\alpha }{𝜆}}\underset{i}{𝜆}{}_{\mu ,\nu }{}^{}.$$ $`(2.7)`$ The second is its dual $`\widehat{\mathrm{\Gamma }}_{.\mu \nu }^\alpha (=\mathrm{\Gamma }_{.\nu \mu }^\alpha )`$, since (2.7) is non-symmetric. The third one is the symmetric part of (2.7), $`\mathrm{\Gamma }_{.(\mu \nu )}^\alpha `$. The fourth is Christoffel symbol defined using (2.4),(2.5) ( as a result of imposing a metricity condition). The torsion tensor is the skew symmetric part of the affine connection (2.7), i.e. $$\mathrm{\Lambda }_{.\mu \nu }^\alpha \stackrel{\mathrm{def}}{=}\mathrm{\Gamma }_{.\mu \nu }^\alpha \mathrm{\Gamma }_{.\nu \mu }^\alpha .$$ $`(2.8)`$ Another third order tensor (contortion) is defined by, $$\gamma _{.\mu \nu }^\alpha \stackrel{\mathrm{def}}{=}\underset{i}{\overset{\alpha }{𝜆}}\underset{i}{𝜆}{}_{\mu ;\nu }{}^{}.$$ $`(2.9)`$ The semicolon is used to characterize covariant differentiation using Christoffel symbol. The two tensors are related by the formula, $$\gamma _{.\mu \nu }^\alpha =\frac{1}{2}(\mathrm{\Lambda }_{.\mu \nu }^\alpha \mathrm{\Lambda }_{\nu .\mu }^\alpha \mathrm{\Lambda }_{\mu .\nu }^\alpha ).$$ $`(2.10)`$ A basic vector could be obtained by contraction of one of the above third order tensors, i.e. $$C_\mu \stackrel{\mathrm{def}}{=}\mathrm{\Lambda }_{.\mu \alpha }^\alpha =\gamma _{.\mu \alpha }^\alpha .$$ $`(2.11)`$ The curvature tensor of the AP-space is, conventionally, defined by, $$B_{.\mu \nu \sigma }^\alpha \stackrel{\mathrm{def}}{=}\mathrm{\Gamma }_{.\mu \sigma ,\nu }^\alpha \mathrm{\Gamma }_{.\mu \nu ,\sigma }^\alpha +\mathrm{\Gamma }_{ϵ\nu }^\alpha \mathrm{\Gamma }_{.\mu \sigma }^ϵ\mathrm{\Gamma }_{.ϵ\sigma }^\alpha \mathrm{\Gamma }_{.\mu \nu }^ϵ0.$$ $`(2.12)`$ This tensor vanishes identically because of (2.1). The autoparallel paths, of this geometry, are given by the equation, $$\frac{d^2x^\mu }{d\lambda ^2}+\mathrm{\Gamma }_{\alpha \beta }^\mu \frac{dx^\alpha }{d\lambda }\frac{dx^\beta }{d\lambda }=0.$$ $`(2.13)`$ The AP-geometry, in its conventional form, has two main problems concerning applications: The first is the identical vanishing of its curvature tensor and the second is that its path equations (2.13) do not represent any known physical trajectory. These problems will be treated in Section 4. Many authors believe that, because of (2.12), the AP-space is flat . It is shown that AP-spaces are, in general, curved. The problem of curvature in AP-spaces is a problem of definition. In any affinely connected space there is, at least, two methods for defining the curvature tensor. The first method is by replacing Christoffel symbol, in the definition of Riemann-Christoffel tensor, by the affine connection defined in the space concerned. The second method is to define curvature as a measure of non-commutation of tensor differentiation using the affine connection of the space. It is known that, the two methods give identical results in case of Riemannian space. But the situation is different for spaces with non-symmetric connections. The two methods are not identical. The application of the second method in non-symmetric geometries implies a problem. That is, we usually use an arbitrary vector in order to study the non-commutation of tensor differentiation, and the resulting expression will not be free from this vector. Fortunately, this problem is solved in AP-spaces . We can replace the arbitrary vector by the vectors defining the structure of AP-spaces. In this case we can define the following curvature tensors (I am going to call these tensors non-conventional curvature tensors): $$\underset{i}{\overset{\stackrel{\mu }{+}}{𝜆}}{}_{\|\nu \sigma }{}^{}\underset{i}{\overset{\stackrel{\mu }{+}}{𝜆}}{}_{\|\sigma \nu }{}^{}=\underset{i}{\overset{\alpha }{𝜆}}B_{.\alpha \nu \sigma }^\mu ,$$ $`(2.14)`$ $$\underset{i}{\overset{\stackrel{\mu }{}}{𝜆}}{}_{\|\nu \sigma }{}^{}\underset{i}{\overset{\stackrel{\mu }{}}{𝜆}}{}_{\|\sigma \nu }{}^{}=\underset{i}{\overset{\alpha }{𝜆}}L_{.\alpha \nu \sigma }^\mu ,$$ $`(2.15)`$ $$\underset{i}{\overset{\stackrel{\mu }{0}}{𝜆}}{}_{\|\nu \sigma }{}^{}\underset{i}{\overset{\stackrel{\mu }{0}}{𝜆}}{}_{\|\sigma \nu }{}^{}=\underset{i}{\overset{\alpha }{𝜆}}N_{.\alpha \nu \sigma }^\mu ,$$ $`(2.16)`$ $$\underset{i}{\overset{\mu }{𝜆}}{}_{;\nu \sigma }{}^{}\underset{i}{\overset{\mu }{𝜆}}{}_{;\sigma \nu }{}^{}=\underset{i}{\overset{\alpha }{𝜆}}R_{.\alpha \nu \sigma }^\mu ,$$ $`(2.17)`$ here we use the stroke , a (+) sign and (-) sign to characterize absolute differentiation using the connection (2.7) and its dual, respectively. We use the stroke without signs to characterize absolute differentiation using the symmetric part of (2.7), while the semicolon is used to characterize covariant differentiation using the Christoffel symbols. The non-conventional curvature tensors defined by (2.14), (2.15), (2.16) and (2.17) are in general non-vanishing except the first one, which vanishes (because of the AP-condition (2.1)). The non-conventional curvature tensors defined above can be written explicitly in terms of torsion, or contortion via (2.10), i.e. $$B_{.\mu \nu \sigma }^\alpha =R_{.\mu \nu \sigma }^\alpha +Q_{.\mu \nu \sigma }^\alpha 0,$$ $`(2.18)`$ $$L_{.\mu \nu \sigma }^\alpha \stackrel{\mathrm{def}}{=}\mathrm{\Lambda }_{.\stackrel{\mu }{+}\stackrel{\nu }{}\|\sigma }^{\stackrel{\alpha }{+}}\mathrm{\Lambda }_{.\stackrel{\mu }{+}\stackrel{\sigma }{}\|\nu }^{\stackrel{\alpha }{+}}+\mathrm{\Lambda }_{.\mu \nu }^\beta \mathrm{\Lambda }_{.\sigma \beta }^\alpha \mathrm{\Lambda }_{.\mu \sigma }^\beta \mathrm{\Lambda }_{.\nu \beta }^\alpha ,$$ $`(2.19)`$ $$N_{.\mu \nu \sigma }^\alpha \stackrel{\mathrm{def}}{=}\mathrm{\Lambda }_{.\mu \nu \|\sigma }^\alpha \mathrm{\Lambda }_{.\mu \sigma \|\nu }^\alpha +\mathrm{\Lambda }_{.\mu \nu }^\beta \mathrm{\Lambda }_{.\beta \sigma }^\alpha \mathrm{\Lambda }_{.\mu \sigma }^\beta \mathrm{\Lambda }_{.\nu \beta }^\alpha ,$$ $`(2.20)`$ $$Q_{.\mu \nu \sigma }^\alpha \stackrel{\mathrm{def}}{=}\gamma _{.\stackrel{\mu }{+}\stackrel{\nu }{+}\|\sigma }^{\stackrel{\alpha }{+}}\gamma _{.\stackrel{\mu }{+}\stackrel{\sigma }{}\|\nu }^{\stackrel{\alpha }{+}}+\gamma _{.\mu \sigma }^\beta \gamma _{.\beta \nu }^\alpha \gamma _{.\mu \nu }^\beta \gamma _{.\beta \sigma }^\alpha ,$$ $`(2.21)`$ It is clear that the vanishing of the torsion will lead to the vanishing of (2.19), (2.20). Also this will lead to vanishing of (2.21) via (2.10) and consequently the vanishing of $`R_{.\mu \nu \sigma }^\alpha `$ via (2.18). This represents another problem facing field theories written in AP-spaces. Such theories will not have GR limit as the torsion vanishes, if this condition is needed. ### 2.2 Quantum Properties of the AP-Geometry Recently , using the affine connections defined in the AP-space to generalize the Bazanski Lagrangian (1.4), three path equations were discovered in the AP-geometry . These equations can be written in the form: $$\frac{dU^\mu }{dS^{}}+\{_{\alpha \beta }^\mu \}U^\alpha U^\beta =0,$$ $`(2.22)`$ $$\frac{dW^\mu }{dS^0}+\{_{\alpha \beta }^\mu \}W^\alpha W^\beta =\frac{1}{2}\mathrm{\Lambda }_{(\alpha \beta ).}^\mu W^\alpha W^\beta ,$$ $`(2.23)`$ $$\frac{dV^\mu }{dS^+}+\{_{\alpha \beta }^\mu \}V^\alpha V^\beta =\mathrm{\Lambda }_{(\alpha \beta ).}^\mu V^\alpha V^\beta ,$$ $`(2.24)`$ where $`S^{}`$,$`S^0`$and $`S^+`$ are the parameters varying along the corresponding curves whose tangents are $`J^\alpha `$,$`W^\alpha `$ and $`V^\alpha `$, respectively. We can write the new set of the path equations, obtained in this geometry, in the following form: $$\frac{dB^\mu }{d\widehat{S}}+a\left\{\genfrac{}{}{0pt}{}{\mu }{\alpha \beta }\right\}B^\alpha B^\beta =b\mathrm{\Lambda }_{(\alpha \beta )}^{..\mu }B^\alpha B^\beta ,$$ $`(2.25)`$ where $`a`$, $`b`$ are the numerical coefficients of the Christoffel symbol term and of the torsion term, respectively. Thus we can construct the following table. Table II: Numerical Coefficients of The Path Equation in AP-Geometry \| Affine Connection Used \| Coefficient $`a`$ \| Coefficient $`b`$ \| \| --- \| --- \| --- \| \| $`\widehat{\mathrm{\Gamma }}_{.\mu \nu }^\alpha `$ \| 1 \| 0 \| \| $`\mathrm{\Gamma }_{.(\mu \nu )}^\alpha `$ \| 1 \| $`\frac{1}{2}`$ \| \| $`\mathrm{\Gamma }_{.\mu \nu }^\alpha `$ \| 1 \| 1 \| The first column in this table contains the affine connections used to generalize the Bazanski Lagrangian. The set of equations (2.22), (2.23)and (2.24) possesses some interesting features: 1. It gives the effect of the torsion on the curves (paths)of the geometry. 2. This set is irreducible i.e. no one of these equations can be reduced to the other unless the torsion vanishes. This vanishing will lead to flat space (in view of the definitions (2.18-21)), which is not suitable for applications. 3. The coefficient of the torsion term jumps by a step of one-half from one equation to the next (as clear from Table II). The last feature is tempting to conclude that: $`\underset{¯}{\mathrm{"}pathsinAPgeometryarenaturallyquantized\mathrm{"}}`$. ### 2.3 Einstein Non-symmetric Geometry Einstein generalized the Riemannian geometry by dropping the symmetry conditions imposed on the metric tensor and on the affine connection . In this geometry the non-symmetric metric tensor is given by: $$g_{\mu \nu }\stackrel{def.}{=}h_{\mu \nu }+f_{\mu \nu },$$ $`(2.26)`$ where, $$h_{\mu \nu }\stackrel{def.}{=}\frac{1}{2}(g_{\mu \nu }+g_{\mu \nu }),$$ $$f_{\mu \nu }\stackrel{def.}{=}\frac{1}{2}(g_{\mu \nu }g_{\mu \nu }).$$ Since the connection of the geometry, $`U_{.\mu \nu }^\alpha `$, is assumed to be non-symmetric, one can define the following 3-types of covariant derivatives: $$A_{+\|\|\nu }^\mu \stackrel{def.}{=}A_{,\nu }^\mu +A^\alpha U_{.\alpha \nu }^\mu ,$$ $`(2.27)`$ $$A_{\|\|\nu }^\mu \stackrel{def.}{=}A_{,\nu }^\mu +A^\alpha U_{.\nu \alpha }^\mu ,$$ $`(2.28)`$ $$A_{0\|\|\nu }^\mu \stackrel{def.}{=}A_{,\nu }^\mu +A^\alpha U_{.(\alpha \nu )}^\mu ,$$ $`(2.29)`$ where $`A^\mu `$ is any arbitrary vector. Now the connection $`U_{.\mu \nu }^\alpha `$ is defined such that , $$g_{\stackrel{\mu }{+}\stackrel{\nu }{}\|\|\sigma }=0,$$ $`(2.30)`$ $$i.e.g_{\mu \nu ,\sigma }=g_{\mu \alpha }U_{.\sigma \nu }^\alpha +g_{\alpha \nu }U_{.\mu \sigma }^\alpha .$$ $`(2.31)`$ The non-symmetric connection can be written in the the following form: $$U_{.\mu \nu }^\alpha \stackrel{def.}{=}U_{.(\mu \nu )}^\alpha +U_{.[\mu \nu ]}^\alpha =\left\{\genfrac{}{}{0pt}{}{\alpha }{\mu \nu }\right\}+K_{.\mu \nu }^\alpha ,$$ $`(2.32)`$ where, $$U_{.(\mu \nu )}^\alpha \stackrel{def.}{=}\frac{1}{2}(U_{.\mu \nu }^\alpha +U_{.\nu \mu }^\alpha ),$$ $`(2.33)`$ $$U_{.[\mu \nu ]}^\alpha \stackrel{def.}{=}\frac{1}{2}(U_{.\mu \nu }^\alpha U_{.\nu \mu }^\alpha )=K_{.[\mu \nu ]}^\alpha =\frac{1}{2}S_{.\mu \nu }^\alpha ,$$ $`(2.34)`$ where $`S_{.\mu \nu }^\alpha `$ is a third order tensor representing the torsion of the Einstein non-symmetric (ENS) geometry. The contravariant metric tensor $`g^{\mu \nu }`$ is defied such that : $$g^{\mu \alpha }g_{\nu \alpha }=g^{\alpha \mu }g_{\alpha \nu }=\delta _\nu ^\mu .$$ $`(2.35)`$ The tensor derivatives (2.27), (2.28) and (2.29) are connected to the parameter derivatives by the relations : $$\frac{A^\mu }{\tau ^{}}=A_{\|\|\alpha }^\mu \stackrel{~}{J}^\alpha ,$$ $`(2.36)`$ $$\frac{A^\mu }{\tau ^0}=A_{0\|\|\alpha }^\mu \stackrel{~}{W}^\alpha ,$$ $`(2.37)`$ $$\frac{A^\mu }{\tau ^+}=A_{+\|\|\alpha }^\mu \stackrel{~}{V}^\alpha ,$$ $`(2.38)`$ where $`\stackrel{~}{J}^\mu `$,$`\stackrel{~}{W}^\mu `$ and $`\stackrel{~}{V}^\mu `$ are tangents to the paths whose evolution parameters are $`\tau ^{}`$,$`\tau ^0`$ and $`\tau ^+`$, respectively. ### 2.4 Quantum Properties of ENS-Geometry Applying the Bazanski approach to the Lagrangian functions: $$\mathrm{\Xi }^{}=g_{\mu \alpha }\stackrel{~}{J}^\mu \frac{\mathrm{\Psi }^\alpha }{\tau ^{}},$$ $`(2.39)`$ $$\mathrm{\Xi }^0=g_{\mu \alpha }\stackrel{~}{W}^\mu \frac{\mathrm{\Theta }^\alpha }{\tau ^0},$$ $`(2.40)`$ $$\mathrm{\Xi }^+=g_{\mu \alpha }\stackrel{~}{V}^\mu \frac{\mathrm{\Phi }^\alpha }{\tau ^+},$$ $`(2.41)`$ where $`\mathrm{\Psi }^\alpha ,\mathrm{\Theta }^\alpha `$ and $`\mathrm{\Phi }^\alpha `$ are the deviation vectors, we get the following set path equations respectively, $$\frac{d\stackrel{~}{J}^\alpha }{d\tau ^{}}+\left\{\genfrac{}{}{0pt}{}{\alpha }{\mu \nu }\right\}\stackrel{~}{J}^\mu \stackrel{~}{J}^\nu =K_{.\mu \nu }^\alpha \stackrel{~}{J}^\mu \stackrel{~}{J}^\nu ,$$ $`(2.42)`$ $$\frac{d\stackrel{~}{W}^\alpha }{d\tau ^0}+\left\{\genfrac{}{}{0pt}{}{\alpha }{\mu \nu }\right\}\stackrel{~}{W}^\mu \stackrel{~}{W}^\nu =\frac{1}{2}g^{\alpha \sigma }g_{\mu \rho }S_{.\nu \sigma }^\rho \stackrel{~}{W}^\mu \stackrel{~}{W}^\nu K_{.\mu \nu }^\alpha \stackrel{~}{W}^\mu \stackrel{~}{W}^\nu ,$$ $`(2.43)`$ $$\frac{d\stackrel{~}{V}^\alpha }{d\tau ^+}+\left\{\genfrac{}{}{0pt}{}{\alpha }{\mu \nu }\right\}\stackrel{~}{V}^\mu \stackrel{~}{V}^\nu =g^{\alpha \sigma }g_{\mu \rho }S_{.\nu \sigma }^\rho \stackrel{~}{V}^\mu \stackrel{~}{V}^\nu K_{.\mu \nu }^\alpha \stackrel{~}{V}^\mu \stackrel{~}{V}^\nu .$$ $`(2.44)`$ This set of equations can be written in the following general form: $$\frac{dC^\alpha }{d\tau }+a\left\{\genfrac{}{}{0pt}{}{\alpha }{\mu \nu }\right\}C^\mu C^\nu =bg^{\alpha \sigma }g_{\mu \rho }S_{.\nu \sigma }^\rho C^\mu C^\nu cK_{.\mu \nu }^\alpha C^\mu C^\nu .$$ $`(2.45)`$ where $`a`$, $`b`$ and $`c`$ are the numerical coefficient of the Christoffel symbol, torsion and K-terms, respectively. Thus, we can construct the following table: Table III: Coefficients of The Path Equations in ENS-Geometry \| Affine Connection used \| Coefficient $`a`$ \| Coefficient $`b`$ \| Coefficient $`c`$ \| \| --- \| --- \| --- \| --- \| \| $`\widehat{U}_{.\mu \nu }^\alpha `$ \| 1 \| 0 \| 1 \| \| $`U_{.(\mu \nu )}^\alpha `$ \| 1 \| $`\frac{1}{2}`$ \| 1 \| \| $`U_{.\mu \nu }^\alpha `$ \| 1 \| 1 \| 1 \| The first column in this table contains the affine connections used to generalize the Bazanski Lagrangian. From this table, it is clear that, the jumping coefficient of the torsion term (column 3) has the same values obtained in the case of the AP-geometry (Table II, column 3). So, one can draw a similar conclusion given in Subsection 2.2: $`\underset{¯}{\mathrm{"}PathsinENSgeometryarenaturallyquantized\mathrm{"}}`$ ## 3 Features of Quantum Roots (i) We consider the jump of the coefficient of the torsion term in the path equations of Subsections 2.2 and 2.4, by a step of one-half, as quantum roots emerging in non-symmetric geometries. Such path equations, are usually used to represent trajectories of test particles, in the context of the scheme of geometerization of physics. So, if such trajectories do exist in nature, then one can conclude that space-time is quantized and the geometry describing nature should be non-symmetric. (ii) The quantum properties shown in Tables II and III, are properties built in the examined geometries. In other words, these properties are intrinsic properties characterizing the type of geometry used. The properties mentioned above are not consequences of applying any known quantization schemes. (iii) In the scheme performed to discover these properties, certain Lagrangian functions are used. Such functions contain, in their structure, covariant derivatives, in which certain affine connections are used. The quantum properties discovered are closely related to such connections. It is well known that, in any non-symmetric geometry, one can define more affine connections by adding any third order tensor to any affine connection already defined in the geometry. If we do so, in the geometries examined in Section 2, one would not get any values (for the coefficients given in Tables II, III) different from those listed in the two tables. As a check one can try the connection, $$\mathrm{\Omega }_{.\mu \nu }^\alpha \stackrel{\mathrm{def}}{=}\{_{\mu \nu }^\alpha \}+\mathrm{\Lambda }_{.\mu \nu }^\alpha ,$$ $`(3.1)`$ defined in the AP-geometry. (iv) As stated above, the quantum properties discovered are closely connected to the affine connection, or more strictly, to its skew pare, the torsion tensor. The coefficients of Christoffel symbol term (the second column of Tables II and III) are the same for all paths. Also, the coefficient of the symmetric part of the tensor $`K_{\mu \nu }^a`$ has no such jumping properties (last column of Table III). ## 4 Parameterization and Diffusion of Quantum Roots It is now obvious that the quantum roots discovered in non-symmetric geometries depend mainly on the existence of non-symmetric connections admitted by such geometries. Furthermore, these roots, explicitly, appeared first in the path equations and not in other geometric entity. In order to extend these roots to the whole geometry, we are going to reconstruct the geometry using a general affine connection. This connection is defined as a linear combination of the connections, already, admitted by the geometry. The combination is carried out using certain parameters. The general expression obtained may not represent an affine connection, in a conventional sense. In other words, it might not be transformed as an affine connection, under the group of general coordinate transformation, unless certain conditions are imposed on the values of the parameters used. The version of the geometry obtained in this way is a parameterized version. In the case of the AP-geometry, using the affine connections mentioned in Subsection 2.1 and carrying out the parameterization scheme mentioned above, the following results are obtained : Combining linearly the above mentioned connections we get, after some reductions, the following parameterized expression, $$_{.\alpha \beta }^\mu =(a+b)\{_{\alpha \beta }^\mu \}+b\gamma _{.\alpha \beta }^\mu $$ $`(4.1a)`$ where $`a`$ and $`b`$ are parameters. As a consequence of imposing a metricity condition, using (4.1a), we get $$a+b=1.$$ $`(4.1b)`$ So, expression (4.1a) will reduce to, $$_{.\alpha \beta }^\mu =\{_{\alpha \beta }^\mu \}+b\gamma _{.\alpha \beta }^\mu ,$$ $`(4.2)`$ which is a general parameterized affine connection. Using (4.2) to generalize the Bazanski Lagrangian (1.4), we get $$\frac{dZ^\mu }{d\tau }+\{_{\alpha \beta }^\mu \}Z^\alpha Z^\beta =b\mathrm{\Lambda }_{(\alpha \beta ).}^\mu Z^\alpha Z^\beta ,$$ $`(4.3)`$ where $`\tau `$ is a parameter varying along the path and $`Z^\mu `$ is the tangent to the path. All curvature tensors defined in this parameterized version of geometry, are non-vanishing. For example if we redefine the curvature (2.12) using the connection (4.2) we get $$B_{}^{}{}_{.\mu \nu \sigma }{}^{\alpha }=R_{.\mu \nu \sigma }^\alpha +b\widehat{Q}_{.\mu \nu \sigma }^\alpha .$$ $`(4.4)`$ where $`R_{.\mu \nu \sigma }^\alpha `$ is Riemann-Christoffel curvature tensor and $`Q_{.\mu \nu \sigma }^\alpha `$ is defined by, $$\widehat{Q}_{.\mu \nu \sigma }^\alpha \stackrel{\mathrm{def}}{=}\gamma _{.\stackrel{\mu }{+}\stackrel{\nu }{+}\|\sigma }^{\stackrel{\alpha }{+}}\gamma _{.\stackrel{\mu }{+}\stackrel{\sigma }{}\|\nu }^{\stackrel{\alpha }{+}}+b(\gamma _{.\mu \sigma }^\beta \gamma _{.\beta \nu }^\alpha \gamma _{.\mu \nu }^\beta \gamma _{.\beta \sigma }^\alpha ).$$ $`(4.5)`$ This tensor is, in general non-vanishing although the corresponding one (2.18) vanishes identically in the old version of the geometry. The torsion and the basic vector of AP-geometry are also parameterized and defined by, $$\mathrm{\Lambda }_{}^{}{}_{.\mu \nu }{}^{\alpha }\stackrel{\mathrm{def}}{=}_{.\mu \nu }^\alpha _{.\nu \mu }^\alpha =b\mathrm{\Lambda }_{.\mu \nu }^\alpha ,$$ $`(4.6)`$ $$C_\mu ^{}\stackrel{\mathrm{def}}{=}\mathrm{\Lambda }_{.\mu \alpha }^\alpha =b\mathrm{\Lambda }_{.\mu \alpha }^\alpha .$$ $`(4.7)`$ The tangent of the new path (4.3) can be written in the form, $$Z^\mu =U^\mu +b\zeta ^\mu ,$$ $`(4.8)`$ where $`U^\mu `$ is the tangent vector of the geodesic of metric and the vector $`\zeta ^\mu `$ represents a deviation from geodesic. The affine parameter $`(\tau )`$ varying along (4.3) can be related to that varying along the geodesic $`(s)`$ by the relation , $$s=\tau (1+bU^\mu \zeta _\mu ).$$ $`(4.9)`$ For physical reasons , the parameter $`b`$ is suggested to take the form $$b=\frac{n}{2}\alpha \gamma ,$$ $`(4.10)`$ where $`n`$ is a natural number, $`\alpha `$ is the fine structure constant and $`\gamma `$ is a dimensionless parameter of order unity.The presence of $`\frac{n}{2}`$ in the parameter $`b`$ will preserve the jumping step appeared in Table II. We are going to call (4.3) the ”Quantum Path Equations”. The torsion term, on the R.H.S. of (4.3), is suggested to represent a type of interaction between the torsion of the background gravitational field and the quantum spin of the moving test particle, Spin-Gravity Interaction. We are going to take $`n=0,1,2,3,\mathrm{}.`$ for particles with spin $`0,\frac{1}{2},1,\frac{3}{2},\mathrm{}.`$, respectively. For slowly rotating macroscopic objects, we are going to take $`n=0`$. ## 5 Quantum Paths and Their Linearization The path equation (4.3) can be used as an equation of motion for any field theory, constructed in the AP-geometry, provided that the theory has good Newtonian limits. In such theories, (e.g. , , ), the tetrad vectors $`\underset{i}{𝜆}_\mu `$ are considered as field variables. So, if we write, $$\underset{i}{𝜆}{}_{\mu }{}^{}=\delta _{{}_{i}{}^{}\mu }+ϵh_{{}_{i}{}^{}\mu },$$ $`(5.1)`$ where $`ϵ`$ is a small parameter, $`\delta _{{}_{i}{}^{}\mu }`$ is Kroneckar delta and $`h_{i\mu }`$ represents deviations from flat space, then the weak field condition can be fulfilled by neglecting quantities of the second and higher orders in $`ϵ`$ in the expanded field quantities. For a static field assumption, we are going to assume the vanishing of time derivatives of the field variables. The vector components $`Z^\mu `$ ($`\stackrel{\mathrm{def}}{=}\frac{dx^\mu }{d\tau }`$)will have the values, $$Z^1Z^2Z^3\epsilon ,Z^01\epsilon ,$$ $`(5.2)`$ where $`\epsilon `$ is a parameter of the order $`(\frac{v}{c})`$. If we want to add the condition of slowly moving particle to the previous conditions we should neglect quantities of second and higher orders of the parameter $`\epsilon `$. Thus, in expanding the quantities of the path equation (4.3) we are going to neglect quantities of orders $`ϵ^2,\epsilon ^2,ϵ\epsilon `$ and higher, and also time derivatives of the field variable are to be neglected. To the first order of the parameters, the only field quantities that will contribute to the path equation (4.3) are given by , $$\mathrm{\Lambda }_{00}^{..i}=ϵh_{00,i},(i=1,2,3)$$ $`(5.3)`$ $$\{_{\mathrm{0\; 0}}^i\}=\frac{ϵ}{2}Y_{00,i},(i=1,2,3)$$ $`(5.4)`$ where $`Y_{\mu \nu }`$ is defined by, $$g_{\mu \nu }=\eta _{\mu \nu }+ϵY_{\mu \nu },$$ $`g_{\mu \nu }`$ is given by (2.5) and $`\eta _{\mu \nu }`$ is the Minkowski metric tensor . Substituting from (5.3),(5.4) into (4.3) we get, after some manipulations : $$\frac{d^2x^i}{d\tau ^2}=\frac{1}{2}ϵ(1\frac{n}{2}\alpha \gamma )Y_{00,i}Z^0Z^0.$$ $`(5.5)`$ In the present case, the metric of the Riemannian space, associated to AP-space, can be written in the form , $$(\frac{d\tau }{dt})^2=c^2(1+ϵY_{00}).$$ $`(5.6)`$ Substituting from (5.6) into (5.5) we get after some manipulations: $$\frac{d^2x^i}{dt^2}=\frac{c^2}{2}ϵ(1\frac{n}{2}\alpha \gamma )Y_{00,i}(i=1,2,3)$$ which can be written in the form, $$\frac{d^2x^i}{dt^2}=\frac{\mathrm{\Phi }_s}{x^i}(i=1,2,3),$$ $`(5.7)`$ where $$\mathrm{\Phi }_s\stackrel{\mathrm{def}}{=}\frac{c^2}{2}ϵ(1\frac{n}{2}\alpha \gamma )Y_{00}.$$ $`(5.8)`$ Equation (5.7) has the same form as Newton’s equation of motion of a particle in a gravitational field having the potential $`\mathrm{\Phi }_s`$ given by (5.8), which differs from the classical Newtonian potential. In the case of motion of macroscopic particles $`(n=0)`$, we get from (5.8): $$\mathrm{\Phi }_s=\frac{c^2}{2}ϵY_{00}=\mathrm{\Phi }_N$$ $`(5.9)`$ where $`\mathrm{\Phi }_N`$ is the Newtonian gravitational potential obtained from a similar treatment of the geodesic equation. Thus (5.8) can be written in the form, $$\mathrm{\Phi }_s=(1\frac{n}{2}\alpha \gamma )\mathrm{\Phi }_N.$$ $`(5.10)`$ This last expression shows that the gravitational potential felt by the spinning particle is less than that felt by a spinless particle or a macroscopic test particle. In other words, the Newtonian potential is reduced, for spinning particles, by a factor $`(1\frac{n}{2}\alpha \gamma )`$. ## 6 Confirmation and Applications of the Quantum Paths In the context of geometerization of physics, path equations of an appropriate geometry, are used to represent trajectories of test particles. It appears clearly, from the previous section, that in the case of a static weak field and a slowly moving test particle, we get Newtonian motion, provided that the particle is spinless. In the following subsections, we are going to use the quantum path equation (4.3), and its linearization consequences, to study the motion of spinning test particles in gravitational fields. ### 6.1 The COW-Experiment Colella, Overhauser and Werner suggested and carried out experiments concerning the quantum interference of thermal neutrons , , . This type of experiments is known, in the literature, as the COW-experiment. The aim of the experiment is to test the effect of the Earth’s gravitational field on the phase difference between two beams of thermal neutrons, one is more closer to the Earth’s surface than the other. The second version of the COW experiment was carried out by Werner et al.. This version is characterized by a high accuracy in the measurements of the phase shift (1 part in 1000). The measurements show that the experimental results are lower than the theoretical calculations (using the Newtonian gravity) by about 8 parts in 1000. This is a real discrepancy, which may indicate the presence of a type of non-Newtonian effects. Now one can use equation (4.3) to give an interpretation for the discrepancy in the COW-experiment. In fact we are going to use the consequence of equation (4.3) given by equation (5.10) since the following conditions, under which (5.10) is derived, hold: -Thermal neutrons can be considered as $`\underset{¯}{slowly}`$ moving test particles, and -the Earth’s gravitational field can be considered as $`\underset{¯}{weak}`$ and $`\underset{¯}{static}`$. The phase difference $`(\mathrm{\Delta }\mathrm{\Omega })`$ between the two beams of neutrons in the COW-experiment is given by (cf. ), $$(\mathrm{\Delta }\mathrm{\Omega })_N=\frac{1}{\mathrm{}}_{ACD}^{ABD}\mathrm{\Phi }_N𝑑t,$$ (6.1) where ABD and ACD are the trajectories of the upper and lower beams of neutrons, in the interferometer, respectively . The index $`N`$ is used to indicate that (6.1) is obtained using the Newtonian potential $`\mathrm{\Phi }_N`$, and $`\mathrm{}`$ is the Planck’s constant. Since neutrons are spinning particles they will be affected by the torsion of space-time, as suggested. Thus we replace $`\mathrm{\Phi }_N`$ in (6.1) by $`\mathrm{\Phi }_S`$ given by (5.10). In this case (6.1) will take the form : $$(\mathrm{\Delta }\mathrm{\Omega })_S=(1\frac{n}{2}\gamma \alpha )\frac{1}{\mathrm{}}_{ACD}^{ABD}\mathrm{\Phi }_N𝑑t,$$ (6.2) i.e., $$(\mathrm{\Delta }\mathrm{\Omega })_S=(\mathrm{\Delta }\mathrm{\Omega })_N\frac{n}{2}\gamma \alpha (\mathrm{\Delta }\mathrm{\Omega })_N.$$ (6.3) The index $`S`$ is used to indicate that (6.2) is obtained using the potential $`\mathrm{\Phi }_S`$. Taking the value of $`\alpha =\frac{1}{137}`$, $`n=1`$ for spin $`\frac{1}{2}`$\- particles (neutrons), we easily get the following results : (1) the theoretical value of the COW-experiment will decrease by about 4 parts in 1000, if we take $`\gamma =1`$, (2) the theoretical value will coincide with the experimental one if we take $`\gamma =2`$. ### 6.2 SN1987A Carriers of astrophysical information are massless spinning particles. These carriers are photons, neutrinos, and expectedly, gravitons. These three types of carriers are assumed to be emitted from supernovae events. In February, $`23^{rd}`$,1987 a supernova , in the Large Magellanic Cloud, was observed (cf.). Observations of the arrival time of neutrinos, at the Kamiokande detectors, was recorded in February $`23^{rd}`$, 1987, $`7^h35^mUT`$, while the arrival time of photons was on the same day at $`10^h40^mUT`$. The bar of the gravitational waves antennae in Rome and Maryland recorded relatively large pulses, 1.2 seconds earlier than neutrinos (cf., ). Although the three types of particles have different spins, general relativity assumes that they follow the same trajectory (null-geodesic of the metric), since they are all massless. In the context of general relativity, it is well known that the time interval required for a massless particle to traverse a given distance is longer in the presence of gravitational field having the potential $`\mathrm{\Phi }(r)`$. The time delay is given by (cf.) $$\mathrm{\Delta }t_{GR}=const._e^a\mathrm{\Phi }(r)𝑑t$$ (6.4) where $`e`$ and $`a`$ are the emission and arrival times of the carrier, respectively. In SN1987A’s time delay (cf. , , ), $`\mathrm{\Phi }(r)`$ is taken to be the Newtonian potential $`\mathrm{\Phi }_N`$ (spin independent). In this case we can construct a spin-independent model, for the emission times of the carriers. If we assume that $`\mathrm{\Phi }(r)`$ is the spin dependent gravitational potential $`\mathrm{\Phi }_s`$ (5.10), we then get the spin-dependent model. The results of these two models are summarized in table IV. Table IV: Emission Times Given By The Models Particles Emitted Spin Independent Model Spin Dependent Model (Cause) (Null-Geodesic) (Quantum Path) Neutrino (core collapse) 0.0 0.0 Photons (maximum brightness) $`+3^h5^m`$ $`+15^h18^m`$ Gravitons (?) $`1^s.2`$ $`+36^h28^m`$ ## From Table IV, we can conclude that, the two models assume two different scenarios for the emission of carriers of astrophysical information. The spin-independent model shows an indication that neutrinos were emitted due to core collapse, associated with gravitons as a result of sudden change in the space-time symmetry, probably, due to a kicked born neutron star. About three hours later photons were emitted as a result of maximum brightness of the envelope. The spin-dependent model shows that: neutrinos were emitted due to a core collapse, preserving sphericity of the core. After 15 hours photons were emitted due to maximum brightness of the envelope, in agreement with SN theories, then 21 hours later the envelope explodes asymmetrically producing a sudden change of space-time symmetry which causes the emission of gravitational waves. It could be seen that the spin-dependent model is more preferable than the spin-independent model. ### 6.3 The Cosmological Parameters Cosmological information are usually carried by, and extracted from, massless spinning particles, ”carriers of cosmological information”. The photon (spin 1-particle) is a good candidate representing one type of these carriers. Recently, the neutrino (spin $`\frac{1}{2}`$-particle) entered the playground as another type. We expect, in the near future, that a third type of carriers, the graviton (spin 2-particle), to be used for extracting cosmological information. Two factors affect the properties of these carriers. The first is the source of the carrier. The second factor is the trajectory of the carrier, in the cosmic space, from its source to the receiver. The first factor implies the information carried, which reflect the properties of the source. The second factor represents the impact of the cosmic space-time on the properties of the carrier. So, information carried by these particles contain, a part connected to their sources, and another part related to the space-time through which these particles travelled. Cosmological parameters are quantities extracted from the information carried by the above mentioned particles. Consequently, the values of such parameters are certainly affected by the second factor. In the present work, we are going to explore the impact of this factor on these parameters. It is well known that, the red-shift of spectral lines, coming from distant objects, plays an important role in measuring the cosmological parameters. Theoretical calculations of the red-shift, in the context of GR, treats it as a metric phenomena, since the metric of space-time is the first integral of the geodesic equation. But, if the trajectories of test particles, the carriers, are spin-dependent, then the red-shift of spectral lines is no longer a metric phenomena. In this case one should look for an alternative scheme for calculating this quantity. Kermack, McCrea and Whittaker developed two theorems on null-geodesics which were applied to get the standard red-shift of relativistic cosmology, using the following formula, $$\frac{\lambda _o}{\lambda _1}=\frac{{}_{}{}^{1}\eta _{}^{\mu }\rho _\mu }{{}_{}{}^{0}\eta _{}^{\mu }\varpi _\mu },$$ $`(6.5)`$ where $`{}_{}{}^{1}\eta _{}^{\mu }`$ is the transport vector along the null-geodesic $`\mathrm{\Gamma }`$ connecting two observers A and B, evaluated at A, $`{}_{}{}^{0}\eta _{}^{\mu }`$ is the transport vector evaluated at B, $`\rho ^\mu `$ is the unit tangent along the trajectory of A, $`\varpi ^\mu `$ is the unit tangent along the trajectory of B, $`\lambda _1`$ is the wave length of the spectral line as measured at A, $`\lambda _o`$ is the wave length of the spectral line as measured at B, and $`\mathrm{\Gamma }`$ represents the trajectory of a massless particle from A (source) to B (receiver). If the universe is expanding then $`\lambda _o>\lambda _1`$. It can be shown that the two theorems, mentioned above, are applicable to any null-path. So, they can be used for massless spinning particles following the trajectory (4.3). In order to evaluate the red-shift using (6.5) one has to know first the values of the vectors used in this formula. Such vectors are obtained as solution of the spin-dependent path equation (4.3). Robertson constructed two geometric AP-structures for cosmological applications. Using one of these structures, and performing the necessary calculations we get , $$\frac{\lambda _o}{\lambda _1}=(\frac{R_o}{R_1})^{(1\frac{n}{2}\alpha \gamma )}$$ $`(6.6).`$ Now, we define the spin-dependent scale factor as, $$R^{}=R^{(1\frac{n}{2}\alpha \gamma ))}.$$ $`(6.7)`$ Using $`R^{}`$, in place of R in the standard definitions of the cosmological parameters, we can list the resulting spin-dependent parameters in Table V. The second column of this Table, gives the values of the parameters as if they are extracted from massless spinless particles. The values of the parameters extracted from photons should match the values listed in column 4. It is worth of mention that the matter parameter is not affected by the spin-gravity interaction. This is due to its independence of Hubble’s parameter. Table V: Spin-Dependent Cosmological Parameters \| Parameter \| Spin-0 \| Spin-$`\frac{1}{2}`$ (neutrino) \| Spin-1 (photon) \| Spin-2 (graviton) \| \| --- \| --- \| --- \| --- \| --- \| \| Hubble \| $`H_o`$ \| $`(1\frac{\alpha }{2})H_o`$ \| $`(1\alpha )H_o`$ \| $`(12\alpha )H_o`$ \| \| Age \| $`\tau _o`$ \| $`\frac{\tau _o}{(1\frac{\alpha }{2})}`$ \| $`\frac{\tau _o}{(1\alpha )}`$ \| $`\frac{\tau _o}{(12\alpha )}`$ \| \| Acceleration \| $`A_o`$ \| $`(1\frac{\alpha }{2})(A_o\frac{\alpha }{2}H_o)`$ \| $`(1\alpha )(A_o\alpha H_o)`$ \| $`(12\alpha )(A_o2\alpha H_o)`$ \| \| Deceleration \| $`q_o`$ \| $`\frac{(q_o\frac{\alpha }{2H_o})}{(1\frac{\alpha }{2})}`$ \| $`\frac{(q_o\frac{\alpha }{H_o})}{(1\alpha )}`$ \| $`\frac{(q_o\frac{2\alpha }{H_o})}{(12\alpha )}`$ \| There are some evidences for the existence of the spin-gravity interaction on the laboratory scale (the results of the COW-experiment), and on the galactic scale (the data of SN1987A). Now, to verify the existence of this interaction on the cosmological scale, observations of one parameter at least, using two different types of carriers, are needed. For example, if we observe neutrinos and photons to get Hubble’s parameter, a discrepancy of order 0.001 would be expected, if this interaction exists on the cosmological scale. ## 7 General Discussion and Concluding Remarks In the present work, it is shown that, starting within the geometerization philosophy, some quantum properties appeared very naturally in the structure of two types of non-symmetric geometries (see the third column of Tables II and III). These properties emerged without imposing any known quantization schemes on the geometry. The properties characterize the torsion term of two new sets of path equations discovered in each geometry, (2.25) and (2.45). The natural appearance of such properties can be considered as quantum roots built in non-symmetric geometry. It is shown that these roots could be extended and diffuse in the whole geometry, using a certain parameterization scheme, suggested in Section 4. This scheme, applied to the AP-geometry, could be applied with some efforts to the ENS-geometry. The application of the parameterization scheme, not only diffuses the quantum properties in the whole geometric structure, but also solves the two main problems of the AP-geometry, mentioned in Subsection 2.1. We can summarize the main advantages of this scheme in the following points: 1. As stated above, it extends the quantum roots, appeared in the path equations, of a non-symmetric geometry, to other geometric entities. 2. It solves completely the curvature problems (the identical vanishing of the curvature (2.12)) , mentioned in Subsection 2.1, by defining a general parameterized non-vanishing curvature tensor (4.4). 3. From the application point of view and depending on the curvature (4.4), field theories written in the parameterized AP-geometry do not need the condition for a vanishing torsion (which leads to a flat AP-space via (2.18-21)), in order to get a correct GR-limit. In other words, to switch-off the effect of the torsion, in such theories, we only take the parameter $`b=0`$. 4. It solves the second problem of conventional AP-geometry, i.e. the non-physical applicability of the path equations (2.13). The new quantum paths (4.3) could be used for physical applications, as shown in Section 6. 5. The parameterized absolute parallelism (PAP) geometry is more general than both the Riemannian and conventional AP-geometries. It could account for both geometries as two limiting cases. These limits can be obtained using (4.1b). The first limit $`a=0b=1`$, which corresponds to the conventional AP-geometry. The second is $`a=1b=0`$, which corresponds to the Riemannian geometry. Figure 1 , is a schematic diagram giving the complete spectrum of geometries admitted by the PAP-geometry. > Figure 1: Quantum Properties of PAP-Geometry. This Figure is plotted using equations (4.4) and (4.6). The new quantum paths (4.3)can be reduced to the geodesic equation of Riemannian geometry (or null-geodesic upon reparameterization), upon setting $`b=0`$, which switch-off the effect of the torsion term. In this case the equation can account for classical mechanics and relativistic mechanics. But if $`b0`$, then the torsion of the background gravitational field will interact with some of the properties of the moving particle. Recalling that the parameter $`b`$ jumps by steps of one-half, (4.10), then one can conclude that the property of the test particle, by which it interacts with the torsion, is its quantum spin. For this reason, the torsion term in (4.3) is suggested to represent ”Spin-Gravity Interaction”. The linearization carried out in Section 5, shows clearly that this interaction will reduce Newton’s gravitational potential, as obvious from (5.10). This equation shows that, the gravitational potential felt by a spinning particle is less than that felt by a spinless particle, or by a macroscopic object. In other words, one can say that, spinning particles feel the space-time torsion. This is similar to the fact that charged particles feel the electromagnetic potential, while neutral particles do not feel it. The discrepancy, between the experimental results and the theoretical calculations (using Newtonian gravity), of the COW-experiment gives a good indicator for the existence of spin-gravity interaction, on the laboratory scale. The experimental results are found to be lower than the expected theoretical calculations. This discrepancy can be interpreted, qualitatively, by a decrease in the gravitational potential, of the Earth, felt by neutrons (spin one-half particles). The value of this potential, felt by neutrons, is less than the value given by Newton’s theory. The application of the new quantum path (4.3), Subsection 6.1, gives good, qualitative and quantitative, agreement with the experimental results. Such agreement gives, not only an evidence of the existence of spin-gravity interaction, but also a direct confirmation of equation (4.3). The application of the linearized form of (4.3), in the case of motion of spinning massless particles coming from SN1987A, Subsection 6.2, gives a good model for the emission times of these particles from this supernova (see Table IV). This may indicate the presence of the spin-gravity interaction on the galactic scale. But more efforts are still needed, both to confirm supernovae mechanisms and for observing more supernovae, to give strong confirmation for the existence of this interaction on the astrophysical scale. The full path equation (4.3) is applied in the case of cosmology, Subsection 6.3. It is shown that the values of the cosmological parameters will be affected by the spin-gravity interaction, if it exists on the cosmological scale. The values of these parameters will depend on the spin of the particle, from which cosmological information are extracted. It is suggested that, a cosmological parameter measured using two massless particles, with different spins (e.g. photon and neutrino) may confirm the existence of spin-gravity interaction on the cosmological scale. The sensitivity of the apparatus, or experiment, to be used should be better than $`0.001`$ In view of the present work, I will try to give short probable answers to some of the good questions raised by professor V.Petrov in the closing session of the conference: Q1: What is the appropriate topology/geometry? A1: A non-symmetric geometry. Q2: How many dimensions? A2: So for, in the context of geometerization of physics, we don’t need more that four dimensions. Mass and charge appear as constants of integration. There are some attempts to represent other interactions (e.g. electromagnetism) together with gravity in spaces of four dimensions (cf. ). Q3: What are the experimental/observational signature of quantum-geometrical effects? A3: Concerning the experimental signature, the COW-type experiment is a good media for testing quantum-geometrical effects on the laboratory scale. the discrepancy in the results of this experiment gives a good indicator for the existence of such effects. Concerning the observational signature, more efforts are needed for observing photons and neutrinos (and probably gravitons, in the future), from supernovae events, in order to detect the existence of such effects on the astrophysical and cosmological scales. Finally, I would like to thank the organizing committee and Professor V.Petrov for inviting me to participate in the conference and to give this talk. References Eisenhart, L.P. 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Colella, R., Overhauser, A.W. and Werner, S.A.(1975) Phys. Rev. Lett., 34,1472. Staudenmann, J.L., Werner, S.A., R. Colella, R. and A. W. Overhauser, A.W. (1980) Phys.Rev. A, 21, 1419. Werner, S.A., Kaiser, H., Arif, M. and Clother, R. (1988) Physica B, 151, 22. Greenberger, D.M. (1983) Rev. Mod. Phys., 55, 875. Wanas, M.I., Melek, M. and Kahil, M.E. (2000) Gravit. Cosmol. 6, 319; gr-qc/9812085. Schramm, D.N. and Truran, J.W. (1990) Phys. Rep. 189, 89-126. Weber, J. (1994) Proc. First Edoarndo Amaldi Conf.”On gravitational wave experiment” Ed. E.Coccia et al. World Scientific P.416. De Rujula, A. (1987) Phys. Lett. 60, 176. Krauss, L.M. and Tremaine, S. (1988) Phys. Rev. Lett., 60, 176. Longo, M.J. (1987) Phys. Rev. D, 36, 3276. Longo, M. J. (1988) Phys. Rev. Lett., 60, 173. Wanas, M.I., Melek, M. and Kahil, M.E. (2002) Proc. MG IX, Vol 2, 1100; gr-qc/0306086. Kermack, W.O., McCrea, W.H. and Whittaker, E.T. (1933) Proc. Roy. Soc. Edin., 53, 31. Robertson, H.P. (1932) Ann. Math. Princeton (2), 33, 496. Wanas, M.I. (2002) To appear in the Proc. IAU-Symp.# 201, held in Manchester, August 2000.
warning/0506/cs0506065.html	ar5iv	text	# Strongly Secure Ramp Secret Sharing Schemes for General Access Structures ## 1 Introduction A secret sharing (SS) scheme is a method to encode a secret $`𝑺`$ into $`n`$ shares each of which has no information of $`𝑺`$, but $`𝑺`$ can be decrypted by collecting several shares. For example, a $`(k,n)`$-threshold SS scheme means that any $`k`$ out of $`n`$ shares can decrypt secret $`𝑺`$ although any $`k1`$ or less shares do not leak out any information of $`𝑺`$. The $`(k,n)`$-threshold access structure can be generalized to so-called general access structures which consist of the families of qualified sets and forbidden sets. A qualified set is the subset of shares that can decrypt the secret, but any information does not leak out from any forbidden set. Generally, the efficiency of SS schemes is evaluated by the entropy of each share, and it must hold that $`H(V_i)H(𝑺)`$ where $`H(𝑺)`$ and $`H(V_i)`$ are the entropies of secret $`𝑺`$ and shares $`V_i`$, $`i=1,2,\mathrm{},n`$, respectively . In order to improve the efficiency of SS schemes, ramp SS schemes are proposed, which have a trade-off between security and coding efficiency . For instance, in the $`(k,L,n)`$-threshold ramp SS scheme , we can decrypt $`𝑺`$ from arbitrary $`k`$ or more shares, but no information of $`𝑺`$ can be obtained from any $`kL`$ or less shares. Furthermore, we assume that arbitrary $`k\mathrm{}`$ shares leak out about $`𝑺`$ with equivocation $`(\mathrm{}/L)H(𝑺)`$ for $`\mathrm{}=1,2,\mathrm{},L`$. In the case where $`L=1`$, the $`(k,L,n)`$-threshold SS scheme reduces to the ordinal $`(k,n)`$-threshold ramp SS scheme. Hence, to distinguish ordinal SS schemes with ramp SS schemes, we call ordinal SS schemes perfect SS schemes. For any $`(k,L,n)`$-threshold access structure, we can realize that $`H(V_i)=H(𝑺)/L`$ , and hence, ramp SS schemes are more efficient than perfect SS schemes . Furthermore, ramp schemes with general access structures are studied in . Since non-forbidden sets with $`1\mathrm{}L1`$ in ramp SS schemes are allowed to leak out a part of a secret, it is important to analyze how the secret partially leaks out. For example, if a secret is a personal data that consists of name, address, job, income, bank account, etc., any part of the secret should not leak out explicitly. However, in the case that the security is measured by the conditional entropy, we cannot know whether or not some part of the secret can be decrypted from a non-forbidden set. Hence, Yamamoto introduced the notion of strong and weak ramp SS schemes . A ramp SS scheme is called a strong ramp SS scheme if it does not leak out any part of a secret explicitly from any arbitrarily $`k\mathrm{}`$ shares for $`\mathrm{}=1,2,\mathrm{},L`$. A ramp SS scheme is weak if it is not strong. But, it is not given how to construct strong ramp SS schemes for arbitrary given general access structures although it is known for $`(k,L,n)`$-threshold ramp SS schemes in . In this paper, we discuss strong ramp SS schemes with general access structures. In section 2, we define ramp SS schemes called partially decryptable (PD) ramp SS schemes, in which every non-qualified set with $`k\mathrm{}`$ shares can decrypt explicitly $`(L\mathrm{})/L`$ parts of a secret. Then, we clarify the relation between PD ramp SS schemes and perfect SS schemes with plural secrets. We also point out that $`(k,L,n)`$-ramp SS schemes based on Shamir’s polynomial interpolation method are not always strong. Next, in section 3, we propose how to convert PD ramp SS schemes into strong ramp SS schemes by using a linear transformation, and we clarify that any access structure that can be realized as a weak ramp SS scheme can also be realized as a strong ramp SS scheme. ## 2 Background and Preliminaries Let $`𝑽=\{V_1,V_2,\mathrm{},V_n\}`$ be the set of all shares, and let $`2^𝑽`$ be the family of all the subsets of $`𝑽`$. Denote a secret by an $`L`$-tuple $`𝑺=\{S_1,S_2,\mathrm{},S_L\}`$, and each element of $`𝑺`$ is assumed to be a mutually independent random variable according to the uniform distribution which takes values in a finite field $`𝔽`$. We assume that $`\|𝔽\|`$ is sufficiently large<sup>1</sup><sup>1</sup>1Throughout this paper, a set of shares and a family of share sets are represented by upper case bold-face and calligraphic font letters, respectively. For simplicity of notation, we use $`𝑨`$$`𝑩`$ to represent $`𝑨𝑩`$ for sets $`𝑨`$ and $`𝑩`$, and $`\{V\}`$ is represented as $`V`$. For example, $`𝑨V=𝑨\{V\}`$. Furthermore, let $`𝑨𝑩`$ be a difference set of $`𝑨`$ and $`𝑩`$, and the cardinality of a set $`𝑨`$ is denoted by $`\|𝑨\|`$.. Then, denote by $`H(𝑺)`$ and $`H(𝑨)`$ the entropies of the secret $`𝑺`$ and a set of shares $`𝑨𝑽`$, respectively. For families $`𝒜_{\mathrm{}}2^𝑽`$, $`\mathrm{}=0,1,\mathrm{},L`$, which consist of subsets of $`𝑽`$, we define ramp SS schemes as follows: ###### Definition 1 Let $`𝑺`$ and $`\mathrm{\Gamma }_L=\{𝒜_0,𝒜_1,\mathrm{},𝒜_L\}`$ be a given secret and a given access structure. Then, $`\{𝑺,𝑽,\mathrm{\Gamma }_L\}`$ is called a ramp secret sharing (SS) scheme if every subset $`𝑨𝒜_{\mathrm{}}`$ satisfies the following for $`\mathrm{}=0,1,\mathrm{},L`$. $`H(𝑺\|𝑨)={\displaystyle \frac{L\mathrm{}}{L}}H(𝑺).`$ (1) $`\mathrm{}`$ Equation (1) implies that secret $`𝑺`$ leaks out from any set $`𝑨𝒜_{\mathrm{}}`$ with the amount of $`(\mathrm{}/L)H(𝑺)`$. Especially, $`𝑺`$ can be completely decrypted from any $`𝑨𝒜_L`$, but any $`𝑨𝒜_0`$ leaks out no information of $`𝑺`$. Hence, in the case of $`L=1`$, ramp SS schemes reduce to perfect SS schemes. Without loss of generality, we can assume that $`𝒜_{\mathrm{}}𝒜_{\mathrm{}^{}}`$ holds for $`\mathrm{}\mathrm{}^{}`$. Furthermore, we also assume that $`_{\mathrm{}=0}^L𝒜_{\mathrm{}}=2^𝑽`$. For example, an access structure of a $`(k,L,n)`$-ramp SS scheme can be defined as $`𝒜_0=\{𝑨:0\|𝑨\|kL\}`$, $`𝒜_{\mathrm{}}=\{𝑨:\|𝑨\|=kL+\mathrm{}\}`$ for $`1\mathrm{}L1`$, and $`𝒜_L=\{𝑨:k\|𝑨\|n\}`$. It is shown in that ramp SS schemes with general access structures can be constructed if and only if the following conditions are satisfied. ###### Theorem 2 () A ramp SS scheme with access structure $`\mathrm{\Gamma }_L=\{𝒜_0,𝒜_1,\mathrm{},𝒜_L\}`$ can be constructed if and only if each $`\stackrel{~}{𝒜}_{\mathrm{}}\stackrel{\mathrm{def}}{=}_{k=\mathrm{}}^L𝒜_k,\mathrm{}=1,2,\mathrm{},L`$ satisfies the monotonicity in the following sense: $`𝑨\stackrel{~}{𝒜}_{\mathrm{}}𝑨^{}\stackrel{~}{𝒜}_{\mathrm{}}\text{for all}𝑨^{}𝑨.`$ (2) $`\mathrm{}`$ In the case of $`L=1`$, (2) in Theorem 2 coincides with the necessary and sufficient condition to realize a perfect SS scheme with an access structure $`\mathrm{\Gamma }_1=\{𝒜_0,𝒜_1\}`$, which is proved in From Theorem 2, the minimal access structure $`𝒜_{\mathrm{}}^{}`$, $`\mathrm{}=1,2,\mathrm{},L`$ can be defined as follows: $`𝒜_{\mathrm{}}^{}=\{𝑨𝒜_{\mathrm{}}:𝑨\{V\}𝒜_{\mathrm{}}\text{for any}V𝑨\}.`$ (3) Proof of Theorem 2 (): We will prove only the sufficiency of (2) because the necessity is clear. Let $`𝑺=\{S_1,S_2,\mathrm{},S_L\}`$ be a secret. From , in the case that (2) holds, we can construct a perfect SS scheme for the secret $`S_{\mathrm{}}`$ with the access structure $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{}}\stackrel{\mathrm{def}}{=}\{2^𝑽\stackrel{~}{𝒜}_{\mathrm{}},\stackrel{~}{𝒜}_{\mathrm{}}\}`$ for every $`\mathrm{}=1,2,\mathrm{},L`$. Then, let $`\stackrel{~}{𝑽}_{\mathrm{}}\stackrel{\mathrm{def}}{=}\{V_{\mathrm{},1},V_{\mathrm{},2},\mathrm{},V_{\mathrm{},n}\}`$ be the set of whole shares for such a perfect SS scheme with access structure $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{}}`$ for the secret $`S_{\mathrm{}}`$. Now, we define $`𝑽_i\stackrel{\mathrm{def}}{=}\{V_{1,i},V_{2,i},\mathrm{},V_{L,i}\}`$ by collecting the $`i`$-th share of $`\stackrel{~}{𝑽}_{\mathrm{}}`$, $`\mathrm{}=1,2,\mathrm{},L`$. Then, it is easy to check that the share set $`𝑽=\{𝑽_1,𝑽_2,\mathrm{},𝑽_n\}`$ realizes the ramp SS scheme with access structure $`\mathrm{\Gamma }_L`$ for the secret $`𝑺`$. In this case, we can decrypt $`\{S_1,S_2,\mathrm{},S_{\mathrm{}}\}`$ from a share set $`𝑨\stackrel{~}{𝒜}_{\mathrm{}}`$, although $`𝑨`$ cannot obtain any information of $`\{S_{\mathrm{}},S_{\mathrm{}+1},\mathrm{},S_L\}`$, and hence, (1) is satisfied. $`\mathrm{}`$ In ramp SS schemes, the coding rate of the $`i`$-th share can be defined as $`\rho _i\stackrel{\mathrm{def}}{=}H(V_i)/H(𝑺)`$. To realize efficient ramp SS schemes, each coding rate of a ramp SS scheme should be as small as possible. Furthermore, it is known that $`\rho _i1/L`$ must hold for each $`i=1,2,\mathrm{},n`$ in any ramp SS scheme with $`L`$-level access structure $`\mathrm{\Gamma }_L`$ . From this viewpoint, the ramp SS schemes shown in the proof of Theorem 2 are not efficient. On the contrary, Okada-Kurosawa presented the following example of a ramp SS scheme with a general access structure, which is more efficient than the ramp SS scheme shown in the proof of Theorem 2. ###### Example 3 () Consider the following access structure $`\mathrm{\Gamma }_2^{\mathrm{ex}}`$ for a set of shares $`𝑽=\{V_1,V_2,V_3,`$ $`V_4\}`$. $`𝒜_1^{}`$ $`=`$ $`\{\{V_1,V_4\},\{V_2,V_4\}\},`$ (4) $`𝒜_2^{}`$ $`=`$ $`\{\{V_1,V_2,V_3\}\}.`$ (5) Then, by letting the secret be $`𝑺=\{S_1,S_2\}`$, a ramp SS scheme for the access structure $`\mathrm{\Gamma }_2^{\mathrm{ex}}`$ in (4) and (5) can be realized as $`V_1`$ $`=`$ $`\{R_1,R_3\},`$ (6) $`V_2`$ $`=`$ $`\{R_2,R_4\},`$ (7) $`V_3`$ $`=`$ $`\{R_1+R_4+S_1,R_2+R_3+S_2\},`$ (8) $`V_4`$ $`=`$ $`\{R_1+S_1,R_2+S_1\},`$ (9) where $`R_1,R_2`$ and $`R_3`$ are mutually independent random numbers which take values in the same finite field $`𝔽`$. $`\mathrm{}`$ From Example 3, it is clear that the secret $`S_2`$ can be decrypted from $`\{V_1,V_4\}`$, but any information of $`S_1`$ cannot be obtained from the set. Hence, since $`S_1`$ and $`S_2`$ are mutually independent, it holds that $`H(𝑺\|V_1V_4)=H(S_1)=H(𝑺)/2`$. In this way, if the partial information of the secret can be explicitly decrypted from every non-qualified set of shares, it is easy to calculate the amount of leaked information. Furthermore, we also note that such a ramp SS scheme can be considered as a special case of perfect SS schemes with $`L`$ plural secrets . In SS schemes with plural secrets, we assume that secret information is given by an $`L`$-tuple $`𝑺^{(L)}=\{S^{(1)},S^{(2)},\mathrm{},S^{(L)}\}`$ where $`S^{(\mathrm{})}`$ are mutually independent random variables. Then, an access structure for the secret $`𝑺^{(L)}`$ is given by $`\mathrm{\Gamma }^{(L)}\stackrel{\mathrm{def}}{=}\{𝒜^{(1)},𝒜^{(2)},\mathrm{},𝒜^{(L)}\}`$ where the secret $`S^{(\mathrm{})}`$ can be decrypted from any set in $`𝒜^{(\mathrm{})}2^𝑽`$ for $`\mathrm{}=1,2,\mathrm{},L`$ while no information of $`S^{(\mathrm{})}`$ can be obtained from any set $`𝑨𝒜^{(\mathrm{})}`$. The SS schemes for $`L`$ secrets with an access structure $`\mathrm{\Gamma }^{(L)}`$ can be defined as follows: ###### Definition 4 () <sup>2</sup><sup>2</sup>2 In the definition of SS schemes with plural secrets in , it is assumed that $`S_{\mathrm{}}`$, $`\mathrm{}=1,2,\mathrm{},L`$, are not always mutually independent. But, we can reduce the definition in to Definition 2, in which $`S_{\mathrm{}}`$’s are mutually independent. Let $`\mathrm{\Gamma }^{(L)}=\{𝒜^{(1)},𝒜^{(2)},\mathrm{},𝒜^{(L)}\}`$ be an access structure for $`L`$ secrets denoted by $`𝑺^{(L)}=\{S^{(1)},S^{(2)},\mathrm{},S^{(L)}\}`$. Then, $`\{𝑺^{(L)},𝑽,\mathrm{\Gamma }^{(L)}\}`$ is called a SS scheme with $`L`$ secrets if it satisfies for all $`\mathrm{}=1,2,\mathrm{}L`$ that $`H(S^{(\mathrm{})}\|𝑨)`$ $`=`$ $`0\text{for any}𝑨𝒜^{(\mathrm{})},`$ (10) $`H(S^{(\mathrm{})}\|𝑨^{})`$ $`=`$ $`H(S^{(\mathrm{})})\text{for any}𝑨^{}𝒜^{(\mathrm{})}.`$ (11) $`\mathrm{}`$ From , Definition 2 is equivalent to the following definition. ###### Definition 5 () Let $`\mathrm{\Gamma }^{(L)}=\{𝒜^{(1)},𝒜^{(2)},\mathrm{},𝒜^{(L)}\}`$ be an access structure for $`L`$ secrets denoted by $`𝑺^{(L)}=\{S^{(1)},S^{(2)},\mathrm{},S^{(L)}\}`$. Let $`𝑺^{(𝑨)}𝑺`$ be a subset of the secret that can be decrypted from a share set $`𝑨𝑽`$ according to $`\mathrm{\Gamma }^{(L)}`$, and we define that $`\overline{𝑺^{(𝑨)}}\stackrel{\mathrm{def}}{=}𝑺𝑺^{(𝑨)}`$. Then, $`\{𝑺^{(L)},𝑽,\mathrm{\Gamma }^{(L)}\}`$ is called a SS scheme with plural secrets $`𝑺^{(L)}`$ if it satisfies that $`H\left(𝑺^{(𝑨)}\|𝑨\right)`$ $`=`$ $`0,`$ (12) $`H\left(\overline{𝑺^{(𝑨)}}\|𝑨\right)`$ $`=`$ $`H\left(\overline{𝑺^{(𝑨)}}\right),`$ (13) for all $`𝑨𝑽`$. $`\mathrm{}`$ Based on Definition 5, we define the partially decryptable ramp SS schemes that characterize the ramp SS schemes shown in the proof of Theorem 2 and Example 3. ###### Definition 6 Let $`𝑺=\{S_1,S_2,\mathrm{},S_L\}`$ be secrets for an access structure $`\mathrm{\Gamma }_L=\{𝒜_1,𝒜_2,\mathrm{},𝒜_L\}`$. Then, $`\{𝑺,𝑽,\mathrm{\Gamma }_L\}`$ is called a partially decryptable (PD) ramp SS scheme if there exists a part of the secret information $`𝑺_𝑨𝑺`$ satisfying that $`\|𝑺_𝑨\|`$ $`=`$ $`\mathrm{}`$ (14) $`H(𝑺_𝑨\|𝑨)`$ $`=`$ $`0,`$ (15) $`H\left(\overline{𝑺_𝑨}\|𝑨\right)`$ $`=`$ $`H\left(\overline{𝑺_𝑨}\right),`$ (16) for all $`𝑨𝒜^{(\mathrm{})}`$ where $`\overline{𝑺_𝑨}\stackrel{\mathrm{def}}{=}𝑺𝑺_𝑨`$. $`\mathrm{}`$ From (15) and (16) in Definition 6, it holds that $`H(𝑺\|𝑨)=H(𝑺_𝑨\|\overline{𝑺_𝑨}𝑨)+H(\overline{𝑺_𝑨}\|𝑨)=H(\overline{𝑺_𝑨})`$, and hence, a PD ramp SS scheme satisfies Definition 1. Note that a PD ramp SS scheme can be regarded as a SS scheme with plural secrets. Conversely, if a SS scheme for plural secrets $`𝑺^{(L)}`$ with access structure $`\mathrm{\Gamma }^{(L)}`$ is given, we can construct a corresponding access structure of a PD ramp SS scheme for the secret $`𝑺=\{S_1,S_2,\mathrm{},S_L\}=\{S^{(1)},S^{(2)},\mathrm{},S^{(L)}\}`$ in the following way: Assign each share set $`𝑨𝑽`$ to the family $`𝒜_{\mathrm{}}`$ where $`\mathrm{}`$ is given by $`\mathrm{}=\left\|\{\mathrm{}^{}:𝑨𝒜^{(\mathrm{}^{})}\mathrm{\Gamma }^{(L)}\}\right\|.`$ (17) Then, the tuple of families $`\{𝒜_0,𝒜_1,\mathrm{},𝒜_L\}\stackrel{\mathrm{def}}{=}\mathrm{\Gamma }_L`$ can be regarded as the access structure of the PD ramp SS scheme. The difference between Definition 5 and Definition 6 is summarized as follows: In Definition 5, from a share set $`𝑨𝑽`$, we can decrypt a subset of secrets $`𝑺^{(L)}`$, i.e., $`𝑺^{(𝑨)}`$, according to the access structure $`\mathrm{\Gamma }^{(L)}`$. However, in the PD ramp SS schemes defined in Definition 6, a share set $`𝑨𝒜_{\mathrm{}}`$ decrypts some $`𝑺_𝑨`$ which satisfies (14), i.e., $`𝑺_𝑨`$ is not specified by the access structure $`\mathrm{\Gamma }_L`$. We note that the amount of the leaked information about $`𝑺`$ from a share set $`𝑨𝒜_{\mathrm{}}`$ is $`(\mathrm{}/L)H(𝑺)`$ in PD ramp SS schemes. Hence, in the sense of (1), there is no difference between Definition 1 and Definition 6. That is, both definitions guarantee the same security in the case that $`𝑺`$ is meaningless if some part of $`𝑺`$ is missing. However, if each part of $`𝑺`$ has explicit meaning, PD ramp SS schemes are not secure, and hence, not desirable. To overcome such defects, Yamamoto defined strong ramp SS schemes as follows <sup>3</sup><sup>3</sup>3In , strong ramp SS schemes are defined for $`(k,L,n)`$-threshold ramp access structures.: ###### Definition 7 () Let $`𝑺=\{S_1,S_2,\mathrm{},S_L\}`$ and $`\mathrm{\Gamma }_L`$ be a secret and an access structure, respectively. Then, $`\{\mathrm{\Gamma }_L,𝑽,𝑺\}`$ is called a strong ramp SS scheme if for all $`\mathrm{}=0,1,\mathrm{},L1`$, $`𝑨𝒜_{\mathrm{}}`$ satisfies (1) and $`H(S_{j_1}S_{j_2}\mathrm{}S_{j_L\mathrm{}}\|𝑨)=H(S_{j_1}S_{j_2}\mathrm{}S_{j_L\mathrm{}})\text{for all}\{S_{j_1},S_{j_2},\mathrm{},S_{j_L\mathrm{}}\}𝑺.`$ (18) $`\mathrm{}`$ Definition 7 implies that strong ramp SS schemes do not leak out any part of the secret explicitly from a non-qualified set $`𝑨𝒜_L`$. Now, from this point of view, we review the $`(k,L,n)`$-threshold SS scheme based on Shamir’s interpolation method. ###### Remark 8 We note that the $`(k,L,n)`$-threshold ramp SS scheme, which is an extension of Shamir’s interpolation method , is not always a strong ramp SS scheme. For instance, consider a $`(4,2,n)`$-threshold ramp SS scheme by using the following polynomial of degree $`3`$ over the finite field $`_{17}`$. $`f(x)=S_1+S_2x+R_1x^2+R_2x^3,`$ (19) where $`𝑺=\{S_1,S_2\}`$ is a secret, and $`R_1`$ and $`R_2`$ are independent random numbers. The $`i`$-th share is given by $`V_i=f(i)`$. Then, from a simple calculation of $`V_3,V_6`$ and $`V_{15}`$, we have $`5S_2=7V_3+9V_6+V_{15}.`$ (20) This means that partial information $`S_2`$ can be decrypted completely from shares $`V_3,V_6`$ and $`V_{15}`$. We also note that from share set $`\{V_1,V_2,V_3\}`$, we have $`H(S_{\mathrm{}}\|V_1V_2V_3)=H(S_{\mathrm{}})`$ for $`\mathrm{}=1,2`$, and hence, the ramp SS scheme in this example is neither PD nor strong<sup>4</sup><sup>4</sup>4In , a construction method is discussed for neither PD nor strong ramp SS schemes.. $`\mathrm{}`$ Remark 8 shows that it is difficult to construct strong ramp SS schemes in general. In , it is proposed how to construct strong $`(k,L,n)`$-threshold ramp SS schemes, but it is not known how to construct strong ramp SS schemes for general access structures. Fortunately, PD ramp SS schemes with general access structure $`\mathrm{\Gamma }_L`$ can easily be constructed if $`\mathrm{\Gamma }_L`$ satisfies monotonicity given by (2) in Theorem 2. Furthermore, it is easy to calculate how much information leaks out from each non-qualified set in PD ramp SS schemes. Therefore, we propose a method to construct strong ramp SS schemes with general access structures based on PD ramp SS schemes. ## 3 Strong Ramp Secret Sharing Schemes with General Access Structures In this section, we propose how to construct a strong ramp SS scheme with general access structure $`\mathrm{\Gamma }_L`$ from a given PD ramp SS scheme with the same access structure $`\mathrm{\Gamma }_L`$. Since a PD ramp SS scheme with general access structure $`\mathrm{\Gamma }_L`$ can always be constructed if $`\mathrm{\Gamma }_L`$ satisfies (2) in Theorem 2, we assume that a PD ramp SS scheme with access structure $`\mathrm{\Gamma }_L=\{𝒜_1,𝒜_2,\mathrm{},𝒜_L\}`$ is obtained for a secret $`𝑺=\{S_1,S_2,\mathrm{},S_L\}`$. Denote by $`\varphi _{\mathrm{\Gamma }_L}(𝑺,𝑹)`$ the encoder of such a PD ramp SS scheme with the access structure $`\mathrm{\Gamma }_L`$ for the secret $`𝑺`$ where $`𝑹`$ represents a set of random numbers used in the encoder. Then, we choose publicly an $`L\times L`$ non-singular matrix $`T`$ and define a new encoder $`\phi _{\mathrm{\Gamma }_L}(𝑺^{},𝑹)\stackrel{\mathrm{def}}{=}\varphi _{\mathrm{\Gamma }_L}(𝑺^{}T,𝑹)`$ where $`𝑺^{}=\{S_1^{},S_2^{},\mathrm{},S_L^{}\}`$<sup>5</sup><sup>5</sup>5 Hereafter, for simplicity of notation, we identify the sets $`𝑺=\{S_1,S_2,\mathrm{},S_L\}`$ and $`𝑺^{}=\{S_1^{},S_2^{},\mathrm{},S_L^{}\}`$ with $`L`$-dimensional row vectors $`[S_1S_2\mathrm{}S_L]`$ and $`[S_1^{}S_2^{}\mathrm{}S_L^{}]`$, respectively.. The next theorem gives the necessary and sufficient condition of $`T`$ that realizes a strong ramp SS scheme with the access structure $`\mathrm{\Gamma }_L`$ for secret $`𝑺^{}=\{S_1^{},S_2^{},\mathrm{},S_L^{}\}`$. ###### Theorem 9 Suppose that the encoder $`\varphi _\mathrm{\Gamma }(𝑺,𝑹)`$ of a PD ramp SS scheme with an access structure $`\mathrm{\Gamma }_L`$ for a secret $`𝑺`$ is given. Let $`𝑺_𝑨`$ be the partial information of the secret $`𝑺`$ that can be decrypted explicitly from a share set $`𝑨`$ in the PD ramp SS scheme, and denote by $`𝑰(𝑨)`$ the set of indices of $`𝑺_𝑨`$. Then, we construct a new encoder $`\phi _{\mathrm{\Gamma }_L}(𝑺^{},𝑹)\stackrel{\mathrm{def}}{=}\varphi _{\mathrm{\Gamma }_L}(𝑺^{}T,𝑹)`$ for a new secret $`𝑺^{}=\{S_1^{},S_2^{},\mathrm{},S_L^{}\}`$ by using a publicly opened $`L\times L`$ non-singular matrix $`T`$. Then, the necessary and sufficient condition of $`T`$ to realize a strong ramp SS scheme $`\{𝑺^{},𝑽,\mathrm{\Gamma }_L\}`$ is given by $`\text{rank}\left[T^1\right]_{j_1,j_2,\mathrm{},j_L\mathrm{}}^{\{1,2,\mathrm{},L\}𝑰(𝑨)}=L\mathrm{},`$ (21) for all $`𝑨𝒜_{\mathrm{}}`$, $`\mathrm{}=0,1,\mathrm{},L`$, where $`\left[T^1\right]_{j_1,j_2,\mathrm{},j_u}^{i_1,i_2,\mathrm{},i_u}`$ is the submatrix that consists of the $`i_1`$-th, $`i_2`$-th$`,\mathrm{},i_u`$-th rows, and the $`j_1`$-th, $`j_2`$-th$`,\mathrm{},j_u`$-th columns of $`T^1`$. $`\mathrm{}`$ ###### Remark 10 Theorem 9 implies that any strong ramp SS schemes can be obtained from the corresponding PD ramp SS schemes without loss of coding rates. $`\mathrm{}`$ Proof of Theorem 9: Since the matrix $`T`$ is non-singular, $`𝑺`$ has one to one correspondence with $`𝑺^{}`$. Hence, $`𝑺^{}`$ is also a set of $`L`$ mutually independent random variables according to the same uniform distribution. Therefore, it holds that $`H(𝑺)=H(𝑺^{})=L\mathrm{log}\|𝔽\|`$ where $`𝔽`$ is a finite field in which $`S_{\mathrm{}}`$, $`\mathrm{}=1,2,\mathrm{},L`$ take values. Then, for any $`𝑨𝒜_{\mathrm{}}`$, $`\mathrm{}=1,2,\mathrm{},L`$, where $`\mathrm{\Gamma }_L=\{𝒜_0,𝒜_1,\mathrm{},𝒜_L\}`$ is the access structure of the PD ramp SS scheme, we have $`H(𝑺^{}\|𝑨)`$ $`=`$ $`H(𝑺\|𝑨)={\displaystyle \frac{L\mathrm{}}{L}}H(𝑺)=(L\mathrm{})\mathrm{log}\|𝔽\|={\displaystyle \frac{L\mathrm{}}{L}}H(𝑺^{}).`$ (22) Therefore, (1) holds for secret $`𝑺^{}`$. Next, from (18), we have for any $`\{S_{j_1}^{},S_{j_2}^{},\mathrm{},S_{j_L\mathrm{}}^{}\}𝑺^{}`$ that $`H(S_{j_1}^{}S_{j_2}^{}\mathrm{}S_{j_L\mathrm{}}^{}\|𝑨)`$ $`=`$ $`H\left(𝑺\left[T^1\right]_{j_1,j_2,\mathrm{},j_L\mathrm{}}^{1,2,\mathrm{},L}\|𝑨\right)`$ (23) $`\stackrel{(\mathrm{a})}{=}`$ $`H\left(\overline{𝑺_𝑨}\left[T^1\right]_{j_1,j_2,\mathrm{},j_L\mathrm{}}^{\{1,\mathrm{},L\}𝑰(𝑨)}\|𝑨\right)`$ $`\stackrel{(\mathrm{b})}{=}`$ $`H\left(\overline{𝑺_𝑨}\|𝑨\right)`$ $`\stackrel{(\mathrm{c})}{=}`$ $`H\left(\overline{𝑺_𝑨}\right)=(L\mathrm{})\mathrm{log}\|𝔽\|=H(S_{j_1}^{}S_{j_2}^{}\mathrm{}S_{j_L\mathrm{}}^{}),`$ where equalities (a), (b), and (c) hold because of (15), (21) and (16), respectively. Finally, we note that the necessity of (21) is clear since equality (b) in (23) does not hold if (21) is not satisfied. $`\mathrm{}`$ From the proof of Theorem 9, it is sufficient to choose the matrix $`T`$ satisfying, instead of the condition (21), that every submatrix of $`T^1`$ has the full rank. We note that the Hilbert matrix $`T_H`$ has such a property. Each element of an $`L\times L`$ Hilbert matrix $`T_H=[t_{ij}]_{\genfrac{}{}{0pt}{}{1iL}{1jL}}`$ is given by $`t_{ij}={\displaystyle \frac{1}{x_i+y_j}},`$ (24) where $`x_i`$ and $`y_j`$ must satisfy for all $`i,j\{1,2,\mathrm{},L\}`$ that $`x_i+y_j0.`$ (25) Note that every submatrix of the Hilbert matrix is also a Hilbert matrix, and the determinant of the matrix $`T_H`$ can be calculated as follows: $`detT_H={\displaystyle \frac{{\displaystyle \underset{1i<jL}{}}(x_ix_j){\displaystyle \underset{1i<jL}{}}(y_iy_j)}{{\displaystyle \underset{i=1}{\overset{L}{}}}{\displaystyle \underset{j=1}{\overset{L}{}}}(x_i+y_j)}}.`$ (26) Hence, it is clear that every submatrix of $`T_H`$ is non-singular if and only if $`x_ix_j\text{and}y_iy_j`$ (27) are satisfied for $`ij`$ in addition to (25). Since $`\|𝔽\|`$ is usually assumed to be sufficiently large in ordinal ramp SS schemes, it is easy to choose $`\{x_i\}_{i=1}^L`$ and $`\{y_i\}_{i=1}^L`$ satisfying (25) and (27). Then, from Theorems 2 and 9, the following theorem holds. ###### Theorem 11 A strong ramp SS scheme with access structure $`\mathrm{\Gamma }_L`$ can be constructed if and only if each $`\stackrel{~}{𝒜_{\mathrm{}}}`$, $`\mathrm{}=1,2,\mathrm{},L`$, satisfies the monotonicity given by (2) of Theorem 2. $`\mathrm{}`$ ###### Example 12 Note that matrices satisfying (21) may exist besides the inverse of Hilbert matrices. As an example, in the case of $`L=2`$ and $`\|𝔽\|3`$, we can use the following matrix $`T^{\mathrm{ex}}`$, the inverse of which is not a Hilbert matrix. $`T^{\mathrm{ex}}=\left[\begin{array}{cc}1& 1\\ 1& 1\end{array}\right].`$ (30) By using the matrix $`T^{\mathrm{ex}}`$ in (30), the PD ramp SS scheme given by (6)–(9) in Example 3 can be transformed into a strong ramp SS scheme with access structure $`\mathrm{\Gamma }_2^{\mathrm{ex}}`$ given by (4) and (5) such that $`V_1=\{R_1,R_3\},V_2=\{R_2,R_4\},V_3=\{R_1+R_4+S_1^{}+S_2^{},R_2+R_3+S_1^{}S_2^{}\}`$, and $`V_4=\{R_1+S_1^{}+S_2^{},R_2+S_1^{}+S_2^{}\}`$. It is easy to check that $`𝑽=\{V_1,V_2,V_3,V_4\}`$ realizes a strong ramp SS scheme with access structure $`\mathrm{\Gamma }_2^{\mathrm{ex}}`$ for secret $`𝑺^{}=\{S_1^{},S_2^{}\}`$. We note here that, in the case of the access structure $`\mathrm{\Gamma }_2^{\mathrm{ex}}`$ in Example 3, the minimum size of $`𝔽`$ is $`2`$ in order realize the PD ramp SS schemes for secret $`𝑺`$ , although $`\|𝔽\|3`$ is required to realize a strong ramp SS schemes for $`𝑺^{}`$ if we use the transformation $`T^{\mathrm{ex}}`$ in (30). In this way, the minimum size of $`𝔽`$ to realize strong ramp SS schemes generally becomes larger than that required to realize PD ramp SS schemes. $`\mathrm{}`$ ###### Remark 13 Note that the matrix $`T`$ described in Theorem 9 is the transformation from a PD ramp SS scheme to a corresponding strong ramp SS scheme. However, weak but not PD ramp SS schemes as shown in Remark 8 cannot always be transformed into strong ramp SS schemes by the matrix $`T`$ satisfying (21). For example, consider the $`(3,2,3)`$-threshold ramp SS scheme given by $`V_1=S_1+R,V_2=S_1+S_2+R`$, and $`V_3=R`$, where $`R`$ is a random number . Then, these shares realize a weak but not PD ramp SS scheme. If we transform this ramp SS scheme by using $`𝑺=𝑺^{}T^{\mathrm{ex}}`$ where $`T^{\mathrm{ex}}`$ is given by (30), we have $`V_1=S_1^{}+S_2^{}+R,V_2=2S_1^{}+R`$, and $`V_3=R`$. It is easy to check that $`V_1,V_2`$ and $`V_3`$ do not realize a strong ramp SS scheme for $`𝑺^{}`$. $`\mathrm{}`$
warning/0506/cond-mat0506694.html	ar5iv	text	# Numerical Renormalization Group for Impurity Quantum Phase Transitions: Structure of Critical Fixed Points ## I Introduction Zero-temperature phase transitions in quantum impurity models have recently attracted considerable interest (for reviews see Refs. BV, ; MV, ; affleck1, ). These transitions can be observed in systems where a zero-dimensional boundary with internal degrees of freedom (the impurity) interacts with an extended bath of fermions or bosons. Examples of impurity models with non-trivial phase transitions include extensions of the Kondo model where one or two magnetic impurities couple to fermionic baths BV , the spin-boson model describing a two-level system coupling to a dissipative environment Leggett ; BTV , as well as so-called Bose-Fermi Kondo models for localized spins interacting with both fermionic and bosonic baths. Impurity phase transitions are of relevance for impurities in correlated bulk systems (e.g. superconductors KV ), for multilevel impurities like Fullerene molecules leo , as well as for nanodevices like coupled quantum dots marcus or point contacts under the influence of dissipative noise karyn ; dima . In addition, impurity phase transitions have been argued to describe aspects of so-called local quantum criticality in correlated lattice systems. Here, the framework of dynamical mean-field theory is employed to map, e.g., the Kondo lattice model onto a single-impurity Bose-Fermi Kondo model supplemented by self-consistency conditions, for details see Ref. edmft, . Diverse techniques have been used to investigate impurity phase transitions, ranging from static and dynamic large-$`N`$ calculations withoff , conformal field theory affleck1 , perturbative renormalization group (RG) BV ; KV ; MVLF and the local-moment approach David ; GL to various numerical methods. In particular, significant progress has been made using the numerical renormalization group (NRG) technique, originally developed by Wilson for the Kondo problem Wil75 . The NRG combines numerically exact diagonalization with the idea of the renormalization group, where progressively smaller energy scales are treated in the course of the calculation. NRG calculations are non-perturbative and are able to access arbitrarily small energies and temperatures. Apart from static and dynamic observables, the NRG provides information about the many-body excitation spectrum of the system at every stage of the RG flow. Thus, it allows to identify fixed points through their fingerprints in the level structure. A detailed understanding of the NRG levels is usually possible if the fixed point can be described by non-interacting bosons or fermions – this is the case for most stable fixed points of impurity models, e.g., the strong-coupling (screened) fixed point of a standard Kondo model. Intermediate-coupling fixed points, usually being interacting, have a completely different NRG level structure, i.e., smaller degeneracies and non-equidistant levels. They cannot be cast into a free-particle description, with the remarkable exception of the two-channel Kondo fixed point which is known to have a representation in terms of free Majorana fermions BHZ . In general, the NRG fixed-point spectrum at impurity transitions is fully universal, apart from a non-universal overall prefactor and discretization effects. The purpose of this paper is to demonstrate that a complete understanding of the NRG many-body spectrum of critical fixed points is actually possible, by utilizing renormalized perturbation theory around a non-interacting fixed point. In the soft-gap Anderson model, this approach can be employed near certain values of the bath exponent which can be identified as critical dimensions. Using the knowledge from perturbative RG calculations, which yield the relevant coupling constants being parametrically small near the critical dimension, we can construct the entire quantum critical many-body spectrum from a free-fermion model supplemented by a small perturbation. In other words, we shall perform epsilon-expansions to determine a complete many-body spectrum (instead of certain renormalized couplings or observables). Vice versa, our method allows to identify relevant degrees of freedom and their marginal couplings by carefully analyzing the NRG spectra near critical dimensions of impurity quantum phase transitions. The paper is organized as follows. In Sec. II we give a brief introduction to the physics of the soft-gap Anderson model and its quantum phase transitions. Sec. III summarizes the recent results from perturbative RG for both the soft-gap Anderson and Kondo models. Section IV describes the numerical renormalization group (NRG) approach which is used here to obtain information about the structure of the quantum critical points. The main part of the paper is Sec. V in which we discuss (i) the numerical data for the structure of the quantum critical points and (ii) the analytical description of these interacting fixed points close to the upper (lower) critical dimension $`r=0`$ ($`r=1/2`$). The main conclusions of the paper are summarized in Sec. VI where we also mention other problems for which an analysis of the type presented here might be useful. ## II Soft-Gap Anderson Model The Hamiltonian of the soft-gap Anderson model withoff is given by: $`H`$ $`=`$ $`\epsilon _f{\displaystyle \underset{\sigma }{}}f_\sigma ^{}f_\sigma +Uf_{}^{}f_{}f_{}^{}f_{}`$ (1) $`+{\displaystyle \underset{k\sigma }{}}\epsilon _kc_{k\sigma }^{}c_{k\sigma }+V{\displaystyle \underset{k\sigma }{}}\left(f_\sigma ^{}c_{k\sigma }+c_{k\sigma }^{}f_\sigma \right).`$ This model describes the coupling of electronic degrees of freedom at an impurity site (operators $`f_\sigma ^{()}`$) to a fermionic bath (operators $`c_{k\sigma }^{()}`$) via a hybridization $`V`$. The $`f`$-electrons are subject to a local Coulomb repulsion $`U`$, while the fermionic bath consists of a non-interacting conduction band with dispersion $`\epsilon _k`$. The model eq. (1) has the same form as the single-impurity Anderson model Hewson but for the soft-gap model we require that the hybridization function $`\stackrel{~}{\mathrm{\Delta }}(\omega )=\pi V^2_k\delta (\omega \epsilon _k)`$ has a soft-gap at the Fermi level, $`\stackrel{~}{\mathrm{\Delta }}(\omega )=\mathrm{\Delta }\|\omega \|^r`$, with an exponent $`r>0`$. This translates into a local conduction band density of states $`\rho (\omega )=\rho _0\|\omega \|^r`$ at low energies. In the numerical calculations we used a band where this power law extends over the whole bandwidth $`D`$, i.e., from $`\omega =D/2`$ to $`+D/2`$, and we have $`\rho _0=(2/D)^{r+1}(r+1)/2`$. However, the universal low-temperature physics to be discussed in the following does not depend on the details of the density of states at high energies, and consequently we will use the low-energy prefactor of the density of states, $`\rho _0`$, to represent the dimensionful energy scale of the problem. Assuming a particle-hole symmetric band, the model (1) is particle-hole symmetric for $`ϵ_f=U/2`$. The soft-gap Anderson model (1) with $`0<r<\mathrm{}`$ displays a very rich behaviour, in particular a continuous transition between a local-moment (LM) and a strong-coupling (SC) phase. Figure 1 shows a typical phase diagram for the soft-gap Anderson model. In the particle-hole (p-h) symmetric case (solid line) the critical coupling $`\mathrm{\Delta }_\mathrm{c}`$ diverges at $`r=\frac{1}{2}`$, and no screening occurs for $`r>\frac{1}{2}`$ (Refs. GBI, ; bullapg, ). No divergence occurs for p-h asymmetry (dashed line) GBI . We now briefly describe the properties of the fixed points in the soft-gap Anderson and Kondo models GBI . Due to the power-law conduction band density of states, already the stable LM and SC fixed points show non-trivial behaviour GBI ; bullapg . The LM phase has the properties of a free spin $`\frac{1}{2}`$ with residual entropy $`S_{\mathrm{imp}}=k_B\mathrm{ln}2`$ and low-temperature impurity susceptibility $`\chi _{\mathrm{imp}}=1/(4k_BT)`$, but the leading corrections show $`r`$-dependent power laws. The p-h symmetric SC fixed point has very unusual properties, namely $`S_{\mathrm{imp}}=2rk_B\mathrm{ln}2`$, $`\chi _{\mathrm{imp}}=r/(8k_BT)`$ for $`0<r<\frac{1}{2}`$. In contrast, the p-h asymmetric SC fixed point simply displays a completely screened moment, $`S_{\mathrm{imp}}=T\chi _{\mathrm{imp}}=0`$. The impurity spectral function follows a $`\omega ^r`$ power law at both the LM and the asymmetric SC fixed point, whereas it diverges as $`\omega ^r`$ at the symmetric SC fixed point – this “peak” can be viewed as a generalization of the Kondo resonance in the standard case ($`r=0`$), and scaling of this peak is observed upon approaching the SC-LM phase boundary bullapg ; David . At the critical point, non-trivial behaviour corresponding to a fractional moment can be observed: $`S_{\mathrm{imp}}=k_B𝒞_S(r)`$, $`\chi _{\mathrm{imp}}=𝒞_\chi (r)/(k_BT)`$ with $`𝒞_S`$, $`𝒞_\chi `$ being universal functions of $`r`$ (see Refs. GBI, ; MVRB, ). The spectral functions at the quantum critical points display an $`\omega ^r`$ power law (for $`r<1`$) with a remarkable “pinning” of the critical exponent. ## III Results from perturbative RG The Anderson model (1) is equivalent to a Kondo model when charge fluctuations on the impurity site are negligible. The Hamiltonian for the soft-gap Kondo model can be written as $`H`$ $`=`$ $`J\stackrel{}{S}\stackrel{}{s_0}+{\displaystyle \underset{k\sigma }{}}\epsilon _kc_{k\sigma }^{}c_{k\sigma }`$ (2) where $`\stackrel{}{s}(0)=_{kk^{}\sigma \sigma ^{}}c_{k\sigma }^{}\stackrel{}{\sigma }_{\sigma \sigma ^{}}c_{k^{}\sigma ^{}}/2`$ is the conduction electron spin at the impurity site $`𝐫=0`$, and the conduction electron density of states follows a power law $`\rho (\omega )=\rho _0\|\omega \|^r`$ as above. ### III.1 RG near $`r=0`$ For small values of the density of states exponent $`r`$, the phase transition in the pseudogap Kondo model can be accessed from the weak-coupling limit, using a generalization of Anderson’s poor man’s scaling. Power counting about the local moment fixed point (LM) shows that $`\mathrm{dim}[J]=r`$, i.e., the Kondo coupling is marginal for $`r=0`$. We introduce a renormalized dimensionless Kondo coupling $`j`$ according to $$\rho _0J=\mu ^rj$$ (3) where $`\mu `$ plays the role of a UV cutoff. The flow of the renormalized Kondo coupling $`j`$ is given by the beta function $$\beta (j)=rjj^2+𝒪(j^3).$$ (4) For $`r>0`$ there is a stable fixed point at $`j^{}=0`$ corresponding to the local-moment phase (LM). An unstable fixed point, controlling the transition to the strong-coupling phase, exists at $$j^{}=r,$$ (5) and the critical properties can be determined in a double expansion in $`r`$ and $`j`$ KV . P-h asymmetry is irrelevant, i.e., a potential scattering term $`E`$ scales to zero according to $`\beta (e)=re`$ (where $`\rho _0E=\mu ^re`$), thus the above expansion captures the p-h symmetric critical fixed point (SCR). As the dynamical exponent $`\nu `$, $`1/\nu =r+𝒪(r^2)`$, diverges as $`r0^+`$, $`r=0`$ plays the role of a lower-critical dimension of the transition under consideration. ### III.2 RG near $`r=1/2`$ For $`r`$ near $`1/2`$ the p-h symmetric critical fixed point moves to strong Kondo coupling, and the language of the p-h symmetric Anderson model becomes more appropriate MVLF . First, the conduction electrons are integrated out exactly, yielding a self-energy $`\mathrm{\Sigma }_f=V^2G_{c0}`$ for the $`f`$ electrons, where $`G_{c0}`$ is the bare conduction electron Green’s function at the impurity location. In the low-energy limit the $`f`$ electron propagator is then given by $$G_f(i\omega _n)^1=i\omega _niA_0\mathrm{sgn}(\omega _n)\|\omega _n\|^r$$ (6) where the $`\|\omega _n\|^r`$ self-energy term dominates for $`r<1`$, and the prefactor $`A_0`$ is $$A_0=\frac{\pi \rho _0V^2}{\mathrm{cos}\frac{\pi r}{2}}.$$ (7) Equation (6) describes the physics of a non-interacting resonant level model with a soft-gap density of states. Interestingly, the impurity spin is not fully screened for $`r>0`$, and the residual entropy is $`2r\mathrm{ln}2`$. This precisely corresponds to the symmetric strong-coupling (SC) phase of the soft-gap Anderson and Kondo model GBI . Dimensional analysis, using $`\mathrm{dim}[f]=(1r)/2`$ \[where $`f`$ represents the dressed fermion according to eq. (6)\], now shows that the interaction term $`U`$ of the Anderson model scales as $`\mathrm{dim}[U]=2r1`$, i.e., it is marginal at $`r=1/2`$. This suggests a perturbative expansion in $`U`$ around the SC fixed point. We introduce a dimensionless renormalized on-site interaction $`u`$ via $$U=\mu ^{2r1}A_0^2u.$$ (8) The beta funcion receives the lowest non-trivial contribution at two-loop order and reads MVLF $`\beta (u)=(12r)u{\displaystyle \frac{3(\pi 2\mathrm{ln}4)}{\pi ^2}}u^3+𝒪(u^5).`$ (9) For $`r<1/2`$ a non-interacting stable fixed point is at $`u^{}=0`$ – this is the symmetric strong-coupling fixed point, it becomes unstable for $`r>1/2`$. Additionally, for $`r<1/2`$ there is a pair of critical fixed points (SCR, SCR’) located at $`u_{}^{}{}_{}{}^{2}=\pi ^2(12r)/[3(\pi 2\mathrm{ln}4)]`$, i.e., $$u^{}=\pm 4.22\sqrt{1/2r}.$$ (10) These fixed points describe the transition between an unscreened (spin or charge) moment phase and the symmetric strong-coupling phase MVLF . Summarizing, both (4) and (9) capture the same critical SCR fixed point. This fixed point can be accessed either by an expansion around the weak-coupling LM fixed point, i.e., around the decoupled impurity limit, valid for $`r1`$, or by an expansion around the strong-coupling SC fixed point, i.e., around a non-interacting resonant-level (or Anderson) impurity, and this expansion is valid for $`1/2r1`$. ## IV Numerical Renormalization Group Here we describe the numerical renormalization group method, suitably extended to handle non-constant couplings $`\stackrel{~}{\mathrm{\Delta }}(\omega )`$ (see Refs. GBI, ; bullapg, ; CJ, ). This method allows a non-perturbative calculation of the many-particle spectrum and physical properties in the whole parameter regime of the model eq. (1), in particular in the low-temperature limit, so that the structure of the quantum critical points is accessible, as discussed in Sec. V. A detailed discussion of how the NRG can be applied to the soft-gap Anderson model can be found in Refs. GBI, ; CJ, ; bullapg, . Here we focus on those aspects of the approach necessary to understand how information on the fixed points can be extracted. The NRG is based on a logarithmic discretization of the energy axis, i.e. one introduces a parameter $`\mathrm{\Lambda }`$ and divides the energy axis into intervals $`[\mathrm{\Lambda }^n,\mathrm{\Lambda }^{(n+1)}]`$ and $`[\mathrm{\Lambda }^{(n+1)},\mathrm{\Lambda }^n]`$ for $`n=0,1,\mathrm{}.,\mathrm{}`$ (see Refs. Wil75, ; Kri80, ). With some further manipulations the original model can be mapped onto a semi-infinite chain with the Hamiltonian $`H`$ $`=\epsilon _f{\displaystyle \underset{\sigma }{}}f_\sigma ^{}f_\sigma +Uf_{}^{}f_{}f_{}^{}f_{}`$ $`+`$ $`\sqrt{{\displaystyle \frac{\xi _0}{\pi }}}{\displaystyle \underset{\sigma }{}}\left[f_\sigma ^{}c_{0\sigma }+c_{0\sigma }^{}f_\sigma \right]`$ $`+`$ $`{\displaystyle \underset{\sigma n=0}{\overset{\mathrm{}}{}}}\left[\epsilon _nc_{n\sigma }^{}c_{n\sigma }+t_n\left(c_{n\sigma }^{}c_{n+1\sigma }+c_{n+1\sigma }^{}c_{n\sigma }\right)\right],`$ with $$\xi _0=_1^1d\omega \stackrel{~}{\mathrm{\Delta }}(\omega ).$$ (12) For a p-h symmetric conduction band, all the on-site energies $`\epsilon _n`$ vanish. If, in addition, the power law in $`\stackrel{~}{\mathrm{\Delta }}(\omega )`$ extends up to a hard cut-off $`\omega _c`$ (we set $`\omega _c=1`$), an exact expression for the hopping matrix elements $`t_n`$ can be givenbullapg : $`t_n`$ $`=`$ $`\mathrm{\Lambda }^{n/2}{\displaystyle \frac{r+1}{r+2}}{\displaystyle \frac{1\mathrm{\Lambda }^{(r+2)}}{1\mathrm{\Lambda }^{(r+1)}}}\left[1\mathrm{\Lambda }^{(n+r+1)}\right]`$ (13) $`\times `$ $`\left[1\mathrm{\Lambda }^{(2n+r+1)}\right]^{1/2}\left[1\mathrm{\Lambda }^{(2n+r+3)}\right]^{1/2}`$ for even $`n`$ and $`t_n`$ $`=`$ $`\mathrm{\Lambda }^{(n+r)/2}{\displaystyle \frac{r+1}{r+2}}{\displaystyle \frac{1\mathrm{\Lambda }^{(r+2)}}{1\mathrm{\Lambda }^{(r+1)}}}\left[1\mathrm{\Lambda }^{(n+1)}\right]`$ (14) $`\times `$ $`\left[1\mathrm{\Lambda }^{(2n+r+1)}\right]^{1/2}\left[1\mathrm{\Lambda }^{(2n+r+3)}\right]^{1/2}`$ for odd $`n`$. The semi-infinite chain is solved iteratively by starting from the impurity and successively adding chain sites. As the coupling $`t_n`$ between two adjacent sites $`n`$ and $`n+1`$ decreases as $`\mathrm{\Lambda }^{n/2}`$ for large $`n`$, the low-energy states of the chain with $`n+1`$ sites are generally determined by a comparatively small number $`N_\mathrm{s}`$ of states close to the ground state of the $`n`$-site system. In practice one retains only these $`N_\mathrm{s}`$ states from the $`n`$-site chain to set up the Hilbert space for the $`n+1`$ site chain, thus preventing the usual exponential growth of the Hilbert space as $`n`$ increases. Eventually, after $`n_{\mathrm{NRG}}`$ sites have been included in the calculation, addition of another site will not change significantly the spectrum of many-particle excitations; the spectrum is very close to that of the fixed point, and the calculation may be terminated. In this way, the NRG iteration gives the many-particle energies $`E_N`$ for a sequence of Hamiltonians $`H_N`$ which correspond to the Hamiltonian eq. (LABEL:eq:H\_si) by the replacement $$\underset{\sigma n=0}{\overset{\mathrm{}}{}}\underset{\sigma n=0}{\overset{N1}{}}.$$ (15) An example for the dependence of the lowest lying energy levels on the chain length (the flow diagram) is given in Fig. 2c for the soft-gap Anderson model with $`r=0.4`$, $`D=2`$, $`U/D=10^3`$ and $`\mathrm{\Delta }=0.0075`$; the parameters used for the NRG calculations are $`\mathrm{\Lambda }=2`$ and $`N_\mathrm{s}=300`$. The states are labelled by the quantum numbers $`Q`$ (which characterizes the number of particles measured relative to half-filling), and the total spin, $`S`$ \[solid lines in Fig. 2 are for $`(Q,S)=(1,0)`$, dashed lines for $`(Q,S)=(0,1/2)`$\]. As mentioned above, the energy scale is reduced in each step by a factor $`\mathrm{\Lambda }^{1/2}`$. To allow for a direct comparison of the energies for different chain lengths, it is thus convenient to plot $`\mathrm{\Lambda }^{N/2}E_N`$ instead of the eigenvalues $`E_N`$ of the $`N`$-site chain directly. Note that here and in the following we use the convention that the energies shown in the flow diagrams are proportional to the bandwidth $`D`$. As is apparent from Fig. 2c, the properties of the system in this case do not change further for chain lengths $`n_{\mathrm{NRG}}>120`$. Without going into details now, one can show that the distribution of energy levels for $`N>120`$ in Fig. 2c is characteristic of the SC phase of the model (see Sec. V). If by contrast we choose instead a value of $`\mathrm{\Delta }=0.006`$, we obtain the flow diagram shown in Fig. 2a. Here it is evident that the fixed point level structure is entirely different from the SC solution, and indeed this particular pattern is now characteristic of the LM phase of the model. We can thus conclude, simply from inspection of the two flow diagrams, that the critical $`\mathrm{\Delta }_\mathrm{c}`$ separating the SC and LM phases of the soft-gap Anderson model for the model parameters specified, lies in the interval $`[0.006,0.0075]`$. Tuning the value of $`\mathrm{\Delta }`$ to the critical value $`\mathrm{\Delta }_\mathrm{c}`$, results in the flow diagram of Fig. 2b. Apparently, the structure of the fixed point at $`\mathrm{\Delta }_\mathrm{c}`$ neither coincides with the SC nor with the LM fixed point. It is clear that it cannot be simply constructed from single-particle states as for the SC and LM fixed points. An important observation is that certain degeneracies present in the SC or LM fixed points are lifted at the QCP. As shown in the following section, a further hint on the structure of the QCPs is given by the dependence of their many-particle spectra on the bath exponent $`r`$. ## V Structure of the Quantum Critical Points In Fig. 3, the many-particle spectra of the three fixed points (SC: dot-dashed lines, LM: dashed lines, and QCP: solid lines) of the symmetric soft-gap model are plotted as functions of the exponent $`r`$ (for a similar figure, see Fig. 13 in Ref. GBI, ). The data are shown for an odd number of sites only and we select the lowest lying energy levels for the subspace $`Q=1`$ and $`S=0`$. As usual, the fixed point structure of the strong coupling and local moment phases can be easily constructed from the single-particle states of a free conduction electron chain. This is discussed in more detail later. Let us now turn to the line of quantum critical points. What information can be extracted from Fig. 3 to understand the structure of these fixed points? First we observe that the levels of the quantum critical points, $`E_{N,\mathrm{QCP}}(r)`$, approach the levels of the LM (SC) fixed points in the limit $`r0`$ ($`r1/2`$): $`\underset{r0}{lim}\left\{E_{N,\mathrm{QCP}}(r)\right\}`$ $`=`$ $`\left\{E_{N,\mathrm{LM}}(r=0)\right\},`$ $`\underset{r1/2}{lim}\left\{E_{N,\mathrm{QCP}}(r)\right\}`$ $`=`$ $`\left\{E_{N,\mathrm{SC}}(r=1/2)\right\},`$ (16) where $`\{\mathrm{}\}`$ denotes the whole set of many-particle states. For $`r0`$, each individual many-particle level $`E_{N,\mathrm{QCP}}(r)`$ deviates linearly from the levels of the LM fixed point, while the deviation from the SC levels is proportional to $`\sqrt{1/2r}`$ for $`r1/2`$. This is illustrated in Fig. 4 where we plot a selection of energy differences $`\mathrm{\Delta }E`$ between levels of QCP and SC fixed points close to $`r=1/2`$. The inset shows the values of the exponents obtained from a fit to the data points. For some levels, there are significant deviations from the exponent $`1/2`$. This is because the correct exponent is only obtained in the limit $`r1/2`$ (the QCP levels have been obtained only up to $`r=0.4985`$). Note that the behaviour of the many-particle levels close to $`r=1/2`$ has direct consequences for physical properties at the QCP; the critical exponent of the local susceptibility at the QCP, for example, shows a square-root dependence on $`(1/2r)`$ close to $`r=1/2`$, see Ref. GBI, . In both limits, $`r0`$ and $`r1/2`$, we observe that degeneracies due to the combination of single-particle levels, present at the LM and SC fixed points, are lifted at the quantum critical fixed points as soon as one is moving away from $`r=0`$ and $`r=1/2`$, respectively. This already suggests that the quantum critical point is interacting and cannot be constructed from non-interacting single-particle states. In the remainder of the paper we want to show how to connect this information from NRG to the perturbative RG of Sec. III. We know that the critical fixed point can be accessed via two different epsilon-expansions KV ; MVLF near the two critical dimensions, and, combined with renormalized perturbation theory, these expansions can be used to evaluate various observables near criticality. Here, we will employ this concept to perform renormalized perturbation theory for the entire many-body spectrum at the critical fixed point. To do so, we will start from the many-body spectrum of one of the trivial fixed points, i.e., LM near $`r=0`$ and SC near $`r=1/2`$, and evaluate corrections to it in lowest-order perturbation theory. This will be done within the NRG concept working directly with the discrete many-body spectra corresponding to a finite NRG chain (which is diagonalized numerically). As the relevant energy scale of the spectra decreases as $`\mathrm{\Lambda }^{n/2}`$ along the NRG iteration, the strength of the perturbation has to be scaled as well, as the goal is to capture a fixed point of the NRG method. This scaling of the perturbation follows precisely from its scaling dimension – the perturbation marginal at the value of $`r`$ corresponding to the critical dimension. With the proper scaling, the operator which we use to capture the difference between the free-fermion and critical fixed points becomes exactly marginal \[see eqs. (21) and (41) below\]. ### V.1 Perturbation theory close to $`r=0`$ Let us now describe in detail the analysis of the deviation of the QCP levels from the LM levels close to $`r=0`$ (the case $`r=1/2`$ is discussed in Sec. V.2). An effective description of the LM fixed point is given by a finite chain with the impurity decoupled from the conduction electron part, see Fig. 5. The conduction electron part of the effective Hamiltonian is given by $$H_{\mathrm{c},N}=\underset{\sigma n=0}{\overset{N1}{}}t_n\left(c_{n\sigma }^{}c_{n+1\sigma }+c_{n+1\sigma }^{}c_{n\sigma }\right).$$ (17) As usual, the structure of the fixed point spectra depends on whether the total number of sites is even or odd. To simplify the discussion in the following, we only consider a total odd number of sites (the flow diagrams of Fig. 2 are all calculated for this case). For the LM fixed point, this means that the number of sites, $`N+1`$, of the free conduction electron chain is even, so $`N`$ in eq. (17) is odd. The single-particle spectrum of the free chain with an even number of sites, corresponding to the diagonalized Hamiltonian $$\overline{H}_{\mathrm{c},N}=\underset{\sigma p}{}ϵ_p\xi _{p\sigma }^{}\xi _{p\sigma },$$ (18) is sketched in Fig. 6. (Note that the $`ϵ_p`$ have to be evaluated numerically.) As we assume p-h symmetry, the positions of the single-particle levels are symmetric with respect to $`0`$ with $$ϵ_p=ϵ_p,p=1,3,\mathrm{},N,$$ (19) and $$\underset{p}{}\underset{p=N,p\mathrm{odd}}{\overset{p=N}{}}.$$ (20) Note that an equally spaced spectrum of single-particle levels is only recovered in the limit $`\mathrm{\Lambda }1`$ (see Fig. 6 in Ref. BHZ, ) for the case $`r=0`$. The RG analysis of Sec. III tells us that the critical fixed point is perturbative accessible from the LM one using a Kondo-type coupling as perturbation. We thus focus on the operator $$H_N^{}=\alpha (r)f(N)\stackrel{}{S}_{\mathrm{imp}}\stackrel{}{s}_0,$$ (21) with the goal to calculate the many-body spectrum of the critical fixed point via perturbation theory in $`H_N^{}`$ for small $`r`$. The function $`\alpha (r)`$ contains the fixed-point value of the Kondo-type coupling, and $`f(N)`$ will be chosen such that $`H_N^{}`$ is exactly marginal, i.e., the effect of $`H_N^{}`$ on the many-particle energies decreases as $`\mathrm{\Lambda }^{N/2}`$ which is the same $`N`$ dependence which governs the scaling of the many-particle spectrum itself. The scaling analysis of Sec. III, eqs. (3) and (5), suggests a parametrization of the coupling as $$\alpha (r)=\frac{\mu ^r}{\rho _0}\alpha r,$$ (22) where $`\rho _0`$ is the prefactor in the density of states, and $`\mu `$ is a scale of order of the bandwidth – such a factor is required here to make $`\alpha `$ a dimensionless parameter. Thus, the strength of the perturbation increases linearly with $`r`$ at small $`r`$ (where $`\mu ^r/\rho _0=D+𝒪(r)`$ for a featureless $`\|\omega \|^r`$ density of states). The qualitative influence of the operator $`\stackrel{}{S}_{\mathrm{imp}}\stackrel{}{s}_0`$ on the many-particle states has been discussed in general in Ref. GBI, for finite $`r`$ and in Refs. Wil75, ; Kri80, for $`r=0`$. Whereas an antiferromagnetic exchange coupling is marginally relevant in the gapless case ($`r=0`$), it turns out to be irrelevant for finite $`r`$, see Ref. GBI, . This is of course consistent with the scaling analysis of Sec. III: the operator (21) simply represents a Kondo coupling, with a tree-level scaling dimension of $`\mathrm{dim}[J]=r`$. A detailed analysis of the $`N`$-dependence of the operator $`\stackrel{}{S}_{\mathrm{imp}}\stackrel{}{s}_0`$ shows that it decreases as $`\mathrm{\Lambda }^{Nr/2}\mathrm{\Lambda }^{N/2}=\mathrm{\Lambda }^{N(r+1)/2}`$ with increasing $`N`$. Consequently, we have to choose $$f(N)=\mathrm{\Lambda }^{Nr/2}.$$ (23) This result also directly follows from $`\mathrm{dim}[J]=r`$: As the NRG discretization yields a decrease of the running energy scale of $`\mathrm{\Lambda }^{N/2}`$, the $`\stackrel{}{S}_{\mathrm{imp}}\stackrel{}{s}_0`$ term in $`H_N^{}`$ (21) scales as $`\mathrm{\Lambda }^{Nr/2}`$. The function $`f(N)`$ is now simply chosen to compensate this effect; using eq. (23) the operator $`H_N^{}`$ becomes exactly marginal. Now we turn to a discussion of the many-body spectrum. The relevant ground state of the effective model for the LM fixed point consists of the filled impurity level (with one electron with either spin $``$ or $``$) and all the conduction electron states with $`p<0`$ filled with both $``$ and $``$, as shown in Fig. 6. Let us now focus on excitations with energy $`ϵ_1+ϵ_3`$ measured with respect to the ground state. Figure 7 shows one such excitation; in this case, one electron with spin $``$ is removed from the $`p=3`$ level and one electron with spin $``$ is added to the $`p=1`$ level. The impurity level is assumed to be filled with an electron with spin $``$, so the resulting state has $`Q=0`$ and $`S_z=+1/2`$. In total, there are 32 states with excitation energy $`ϵ_1+ϵ_3`$. These states can be classified using the quantum numbers $`Q`$, $`S`$, and $`S_z`$. Here we consider only the states with quantum numbers $`Q=0`$, $`S=1/2`$, and $`S_z=1/2`$ (with excitation energy $`ϵ_1+ϵ_3`$) which form a four-dimensional subspace. As the state shown in Fig. 7 is not an eigenstate of the total spin $`S`$, we have to form proper linear combinations to obtain a basis for this subspace; this basis can be written in the form: $`\|\psi _1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f_{}^{}\left(\xi _1^{}\xi _3+\xi _1^{}\xi _3\right)\|\psi _0`$ $`\|\psi _2`$ $`=`$ $`\left[{\displaystyle \frac{1}{\sqrt{6}}}f_{}^{}\left(\xi _1^{}\xi _3\xi _1^{}\xi _3\right)+{\displaystyle \frac{2}{\sqrt{6}}}f_{}^{}\xi _1^{}\xi _3\right]\|\psi _0`$ $`\|\psi _3`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f_{}^{}\left(\xi _3^{}\xi _1+\xi _3^{}\xi _1\right)\|\psi _0`$ $`\|\psi _4`$ $`=`$ $`\left[{\displaystyle \frac{1}{\sqrt{6}}}f_{}^{}\left(\xi _3^{}\xi _1\xi _3^{}\xi _1\right)+{\displaystyle \frac{2}{\sqrt{6}}}f_{}^{}\xi _3^{}\xi _1\right]\|\psi _0`$ where the state $`\|\psi _0`$ is given by the product of the ground state of the conduction electron chain and the empty impurity level: $$\|\psi _0=\left[\underset{p<0}{}\xi _p^{}\xi _p^{}\|0_{\mathrm{cond}}\right]\|0_{\mathrm{imp}}.$$ (25) The fourfold degeneracy of the subspace ($`Q=0`$, $`S=1/2`$, $`S_z=1/2`$) of the LM fixed point at energy $`ϵ_1+ϵ_3`$ is partially split for finite $`r`$ in the spectrum of the quantum critical fixed point. Let us now calculate the influence of the perturbation $`H_N^{}`$ on the states $`\|\psi _1,\mathrm{}\|\psi _4`$, concentrating on the splitting of the energy levels up to first order. Degenerate perturbation theory requires the calculation of the matrix $$W_{ij}=\psi _i\|H_N^{}\|\psi _j,i,j=1,\mathrm{}4,$$ (26) and a subsequent calculation of the eigenvalues of $`\left\{W_{ij}\right\}`$ gives the level splitting. Details of the calculation of the matrix elements $`W_{ij}`$ are given in Appendix A. The result is $$\left\{W_{ij}\right\}=\alpha (r)f(N)\left[\begin{array}{cccc}0& \frac{\sqrt{3}}{4}\gamma & 0& 0\\ \frac{\sqrt{3}}{4}\gamma & \frac{1}{2}\beta & 0& 0\\ 0& 0& 0& \frac{\sqrt{3}}{4}\gamma \\ 0& 0& \frac{\sqrt{3}}{4}\gamma & \frac{1}{2}\beta \end{array}\right],$$ (27) with $`\gamma =\left[\|\alpha _{01}\|^2\|\alpha _{03}\|^2\right]`$ and $`\beta =\left[\|\alpha _{01}\|^2+\|\alpha _{03}\|^2\right]`$. The $`N`$-dependence of the coefficients $`\alpha _{0p}`$ (which relate the operators $`c_{0\sigma }`$ and $`\xi _{p\sigma }`$, see eq. (61)) is given by $$\|\alpha _{0p}\|^2\mathrm{\Lambda }^{Nr/2}\mathrm{\Lambda }^{N/2},$$ (28) (see also Sec. III A in Ref. GBI, ). Numerically we find that $`\gamma `$ $`=`$ $`0.1478\mathrm{\Lambda }^{Nr/2}\mathrm{\Lambda }^{N/2}`$ $`\beta `$ $`=`$ $`2.0249\mathrm{\Lambda }^{Nr/2}\mathrm{\Lambda }^{N/2},`$ (29) where the prefactors depend on the exponent $`r`$ and the discretization parameter $`\mathrm{\Lambda }`$ (the quoted values are for $`r=0.01`$ and $`\mathrm{\Lambda }=2.0`$). The matrix $`\left\{W_{ij}\right\}_{r=0.01}`$ then takes the form $`\left\{W_{ij}\right\}_{r=0.01}=\alpha (r=0.01)\mathrm{\Lambda }^{N/2}`$ (34) $`\times `$ $`\left[\begin{array}{cccc}0& 0.064& 0& 0\\ 0.064& 1.013& 0& 0\\ 0& 0& 0& 0.064\\ 0& 0& 0.064& 1.013\end{array}\right].`$ Diagonalization of this matrix gives the first-order corrections to the energy levels $`\mathrm{\Delta }E_1(r=0.01)`$ $`=`$ $`\mathrm{\Delta }E_3(r=0.01)`$ $`=`$ $`\alpha (r=0.01)\mathrm{\Lambda }^{N/2}(1.0615)`$ $`\mathrm{\Delta }E_2(r=0.01)`$ $`=`$ $`\mathrm{\Delta }E_4(r=0.01)`$ (35) $`=`$ $`\alpha (r=0.01)\mathrm{\Lambda }^{N/2}0.0004`$ with $`E_{N,\mathrm{QCP}}(r=0.01,i)`$ $`=`$ $`E_{N,\mathrm{LM}}(r=0.01,i)`$ (36) $`+`$ $`\mathrm{\Delta }E_i(r=0.01),`$ ($`i=1,\mathrm{}4`$). Apparently, the fourfold degeneracy of the subspace ($`Q=0`$, $`S=1/2`$, $`S_z=1/2`$) with energy $`ϵ_1+ϵ_3`$ is split in two levels which are both twofold degenerate. We repeated this analysis for a couple of other subspaces and a list of the resulting matrices $`\left\{W_{ij}\right\}`$ and the energy shifts $`\mathrm{\Delta }E`$ is given in Appendix A. Let us now proceed with the comparison of the perturbative results with the structure of the quantum critical fixed point calculated from the NRG. For our specific choice of the conduction band density of states, the relation (22) yields $`\alpha (r)=\alpha rD`$ for small $`r`$ (where $`\mu ^r1`$). Using the corresponding equations for the energy shifts in Appendix A, we observe that a single parameter $`\alpha `$ must be sufficient to describe the level shifts in all subspaces, provided that the exponent $`r`$ is small enough so that the perturbative calculations are still valid. A numerical fit gives $`\alpha 1.03`$ for $`\mathrm{\Lambda }=2.0`$, (the $`\mathrm{\Lambda }`$-dependence of $`\alpha `$ is discussed later, see Fig. 9). Figure 8 summarizes the NRG results together with the perturbative analysis for exponents $`r`$ close to 0. A flow diagram of the lowest lying energy levels is shown in Fig. 8a for a small value of the exponent, $`r=0.03`$, so that the levels of the QCP only slightly deviate from those of the LM fixed point. As discussed above, the deviation of the QCP levels from the LM levels increases linearly with $`r`$, see Fig. 8b. We indeed get a very good agreement between the perturbative result (crosses) and the NRG-data (lines) for exponents up to $`r0.07`$. The data shown here are for the subspaces ($`Q=0`$, $`S=1/2`$, $`S_z=1/2`$) and energy $`2ϵ_1`$ (the levels at $`E_N\mathrm{\Lambda }^{N/2}1`$, see Appendix A.1) and ($`Q=0`$, $`S=1/2`$, $`S_z=1/2`$) and energy $`ϵ_1+ϵ_3`$ (the levels at $`E_N\mathrm{\Lambda }^{N/2}2`$, see the example discussed in this section). In the NRG, the continuum limit corresponds to the limit $`\mathrm{\Lambda }1`$, but due to the drastically increasing numerical effort upon reducing $`\mathrm{\Lambda }`$, results for the continuum limit have to be obtained via extrapolation of NRG data for $`\mathrm{\Lambda }`$ in, for example, the range $`1.5<\mathrm{\Lambda }<3.0`$. Figure 9 shows the numerical results from the NRG calculation together with a linear fit to the data: $`\alpha (\mathrm{\Lambda })=0.985+0.045(\mathrm{\Lambda }1.0)`$. Taking into account the increasing error bars for smaller values of $`\mathrm{\Lambda }`$, the extrapolated value $`\alpha (\mathrm{\Lambda }1)0.985`$ is in excellent agreement with the result from the perturbative RG calculation, which is directly for the continuum limit and gives $`\alpha =1.0`$. ### V.2 Perturbation theory close to $`r=1/2`$ To describe the deviation of the QCP levels from the SC levels close to $`r=1/2`$, we have to start from an effective description of the SC fixed point. This is given by a finite chain including the impurity site with the Coulomb repulsion $`U=0`$ at the impurity site and a hybridization $`\overline{V}`$ between impurity and the first conduction electron site, see Fig. 10. Note that the SC fixed point can also be described by the limit $`\overline{V}\mathrm{}`$ and finite $`U`$ which means that impurity and first conduction electron site are effectively removed from the chain. This reduces the number of sites of the chain by two and leads to exactly the same level structure as including the impurity with $`U=0`$. However, the description with the impurity included (and $`U=0`$) is more suitable for the following analysis. The corresponding effective Hamiltonian is that of a soft-gap Anderson model on a finite chain with $`N+2`$ sites and $`\epsilon _f=U=0`$ (i.e., a p-h symmetric resonant level model). $$H_{\mathrm{sc},N}=\overline{V}\underset{\sigma }{}\left[f_\sigma ^{}c_{0\sigma }+c_{0\sigma }^{}f_\sigma \right]+H_{\mathrm{c},N},$$ (37) with $`H_{\mathrm{c},N}`$ as in eq. (17). As for the effective description of the LM fixed point, the effective Hamiltonian is that of a free chain. Focussing, as above, on odd values of $`N`$, the total number of sites of this chain, $`N+2`$, is odd. The single-particle spectrum of the free chain with an odd number of sites, corresponding to the diagonalized Hamiltonian $$\overline{H}_{\mathrm{sc},N}=\underset{\sigma l}{}ϵ_l\xi _{l\sigma }^{}\xi _{l\sigma },$$ (38) is sketched in Fig. 11. As we assume p-h symmetry, the positions of the single-particle levels are symmetric with respect to $`0`$ with $$ϵ_0=0,ϵ_l=ϵ_l,l=2,4,\mathrm{},(N+1),$$ (39) and $$\underset{l}{}\underset{l=(N+1),l\mathrm{even}}{\overset{l=N+1}{}}.$$ (40) The ground state of the effective model for the SC fixed point is fourfold degenerate, with all levels with $`l<0`$ filled and the level $`l=0`$ either empty, singly ($``$ or $``$) or doubly occupied. According to Sec. III the proper perturbation to access the critical fixed point from the SC one is an on-site repulsion, thus we choose $$H_N^{}=\beta (r)\overline{f}(N)\left(n_f\frac{1}{2}\right)\left(n_f\frac{1}{2}\right),$$ (41) ($`n_{f\sigma }=f_\sigma ^{}f_\sigma `$) with the strength of the perturbation parametrized as $$\beta (r)=\mu ^{2r1}\rho _0^2\overline{V}^4\beta \sqrt{1/2r}.$$ (42) see Sec. III. Note that $`\rho _0^2(r=1/2)=9/(2D^3)`$ for a featureless power-law density of states with bandwidth $`D`$. The $`N`$ dependence of the operator $`\left(n_f\frac{1}{2}\right)\left(n_f\frac{1}{2}\right)`$ is given by $`\mathrm{\Lambda }^{(r1/2)N}\mathrm{\Lambda }^{N/2}=\mathrm{\Lambda }^{(r1)N}`$, so we have to choose $$\overline{f}(N)=\mathrm{\Lambda }^{(1/2r)N}.$$ (43) This again follows from the scaling analysis of Sec. III: the on-site repulsion has scaling dimension $`\mathrm{dim}[U]=2r1`$. Thus the $`\left(n_f\frac{1}{2}\right)\left(n_f\frac{1}{2}\right)`$ term in $`H_N^{}`$ (41) scales as $`\mathrm{\Lambda }^{N(r1/2)}`$, and $`\overline{f}(N)`$ (43) compensates this behavior to make $`H_N^{}`$ exactly marginal. We continue with analyzing the low-lying many-body levels. Similar as above, we focus on one specific example, these are excitations with energy $`2\epsilon _2`$ measured with respect to the ground state and quantum numbers $`Q=1`$, $`S=0`$, and $`S_z=0`$. This subspace is two-dimensional and the basis is given by $`\|\psi _1`$ $`=`$ $`\xi _0^{}\xi _0^{}\xi _2\xi _2\|\psi _0,`$ $`\|\psi _2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\xi _2^{}\xi _2+\xi _2^{}\xi _2\right)\|\psi _0,`$ (44) with $$\|\psi _0=\underset{l<0}{}\xi _l^{}\xi _l^{}\|0.$$ (45) (Note that in this definition of $`\|\psi _0`$, the $`l=0`$-level is empty.) The two-fold degeneracy of this subspace is lifted for $`r<1/2`$ in the spectrum of the quantum critical points. The matrix $`W_{ij}=\psi _i\|H_N^{}\|\psi _j`$ ($`i,j=1,2`$) is given by $$\left\{W_{ij}\right\}=\beta (r)\overline{f}(N)\|\alpha _{f2}\|^4\left[\begin{array}{cc}22\kappa +\kappa ^2& 2\sqrt{2}\kappa \\ 2\sqrt{2}\kappa & 2+\kappa ^2\end{array}\right],$$ (46) with $`\kappa =\|\alpha _{f0}\|^2/\|\alpha _{f2}\|^2`$. The N-dependence of the coefficients $`\|\alpha _{fl}\|`$ (which relate the operators $`f_\mu `$ and $`\xi _{l\mu }`$, see eq. (85)) is given by $$\|\alpha _{fl}\|^2\mathrm{\Lambda }^{(r1)N/2},$$ (47) Numerically we find that $`\|\alpha _{f2}\|^2`$ $`=`$ $`0.0366(D/\overline{V})^2\mathrm{\Lambda }^{(r1)N/2}`$ $`\|\alpha _{f0}\|^2`$ $`=`$ $`0.0930(D/\overline{V})^2\mathrm{\Lambda }^{(r1)N/2},`$ (48) where the prefactors depend on the exponent $`r`$ and the quoted value is for $`r=0.499`$. The matrix $`\left\{W_{ij}\right\}_{r=0.499}`$ then takes the form $`\left\{W_{ij}\right\}_{r=0.499}`$ $`=`$ $`\beta (r=0.499)(D/\overline{V})^4\mathrm{\Lambda }^{N/2}`$ (51) $`\times `$ $`\left[\begin{array}{cc}0.0044& 0.0094\\ 0.0094& 0.011\end{array}\right],`$ Diagonalization of this matrix gives the first-order corrections to the energy levels $`\mathrm{\Delta }E_1(r=0.499)`$ $`=`$ $`\beta (r=0.499)(D/\overline{V})^4\mathrm{\Lambda }^{N/2}(0.0023)`$ $`\mathrm{\Delta }E_2(r=0.499)`$ $`=`$ $`\beta (r=0.499)(D/\overline{V})^4\mathrm{\Lambda }^{N/2}(0.018)`$ with $`E_{N,\mathrm{QCP}}(r=0.499,i)`$ $`=`$ $`E_{N,\mathrm{SC}}(r=0.499,i)`$ $`+`$ $`\mathrm{\Delta }E_i(r=0.499),`$ (53) ($`i=1,2`$). We repeated this analysis for a couple of other subspaces and a list of the resulting matrices $`\left\{W_{ij}\right\}`$ and the energy shifts $`\mathrm{\Delta }E`$ is given in Appendix B. The comparison of the perturbative results with the numerical results from the NRG calculation is shown in Fig. 12b. As for the case $`r0`$ we observe that a single parameter $`\beta `$ is sufficient to describe the level shifts in all subspaces, provided the exponent $`r`$ is close enough to $`r=1/2`$ so that the perturbative calculations are valid. For $`\mathrm{\Lambda }=2.0`$ we find $`\beta 70`$ and the $`\mathrm{\Lambda }1`$ extrapolation results in $`\beta (\mathrm{\Lambda }1)73.0\pm 5.0`$ (the error bars here are significantly larger as for the extrapolation of the coupling $`\alpha `$). The results from perturbative RG, Sec. III, specifically eqs. (8) and (10), yield $`\beta (r)=\mu ^{2r1}\rho _0^2\overline{V}^4\mathrm{\hspace{0.17em}2}\pi ^2u^{}`$. This gives $`\beta =83.3`$. Similar to Fig. 8 above, we show in Fig. 12a a flow diagram for an exponent very close to $`1/2`$, $`r=0.4985`$, so that the levels of the QCP only slightly deviate from those of the SC levels. As discussed above, this deviation is proportional to $`\sqrt{1/2r}`$, see Fig. 12b. The data shown here are all for subspaces with ($`Q=1`$, $`S=0`$, $`S_z=0`$); the unperturbed energies $`E`$ of these subspaces are: * $`E=0`$: the levels at $`E_N\mathrm{\Lambda }^{N/2}0`$, see App. B.2, * $`E=ϵ_2`$: the levels at $`E_N\mathrm{\Lambda }^{N/2}0.8`$, see App. B.3, * $`E=2ϵ_2`$: the levels at $`E_N\mathrm{\Lambda }^{N/2}1.6`$, see the example discussed in this section, * $`E=ϵ_4`$: the levels at $`E_N\mathrm{\Lambda }^{N/2}1.8`$, see App. B.4, * $`E=3ϵ_2`$: the levels at $`E_N\mathrm{\Lambda }^{N/2}2.4`$. We again find a very good agreement between the perturbative results (crosses) and the NRG data (lines). Thus we can summarize that our renormalized perturbation theory for the NRG many-body spectrum works well in the vicinity of both $`r=0`$ and $`r=1/2`$. In principle, from the many-body spectrum (and suitable matrix elements) all other observables like thermodynamic data and dynamic correlation functions can be calculated. We note that the convergence radius of the epsilon-expansion for the levels seems to be smaller than that of the direct epsilon-expansion for certain observables like critical exponents and impurity susceptibility and entropy, see Ref. MVLF, . ## VI Conclusions Using the quantum phase transitions in the soft-gap Anderson model as an example, we have demonstrated that epsilon-expansion techniques can be used to determine complete many-body spectra at quantum critical points. To this end, we have connected information from standard perturbative RG, which yields information on critical dimensions and parametrically small couplings, and from NRG for the many-body spectra of free-fermion fixed points. Together, these can be used to perform renormalized perturbation theory for many-body spectra of interacting intermediate-coupling fixed points. For the soft-gap Anderson model, which features two lower-critical dimensions at $`r=0`$ and $`r=1/2`$, correspondingly two different approaches can be utilized to capture the same critical fixed point: Near $`r=0`$ a Kondo term has to be added to a free-fermion chain with a decoupled impurity, whereas near $`r=1/2`$ an on-site repulsion is used as a perturbation to the non-interacting Anderson (or resonant-level) model. These perturbations lift the large degeneracies present in the non-interacting spectra, and accurately reproduce the critical spectra determined in NRG calculations at criticality. Vice versa, our method will be useful in situations where the effective low-energy theory for the critical point is not known: a careful analysis of the many-body spectrum near critical dimensions yields information about the scaling dimension and structure of the relevant operators. For instance, a plot similar to Fig. 3 can be calculated for the spin-boson model, using the numerical renormalization group method as in Ref. BTV, . Preliminary results (not shown here) indicate that the many-particle levels of the QCP approach the levels of the delocalized (localized) fixed point in the limit $`s0`$ ($`s1`$), with $`s`$ the exponent of the bath spectral function $`J(\omega )\omega ^s`$. We envision applications of our ideas to more complex impurity models, e.g., with two orbitals or two coupled spins, as well as to non-equilibrium situations treated using NRG fba . ###### Acknowledgements. We thank S. Kehrein, Th. Pruschke, and A. Rosch for discussions and S. Florens, L. Fritz, M. Kirćan, and N. Tong for collaborations on related work. This research was supported by the DFG through SFB 484 (HJL, RB) and the Center for Functional Nanostructures Karlsruhe (MV). MV also acknowledges support from the Helmholtz Virtual Quantum Phase Transitions Institute in Karlsruhe. ## Appendix A Details of the Perturbative Analysis around the Local Moment Fixed Point In this Appendix, we want to give more details for the derivation of the matrix $`W_{ij}`$ eq. (27) which determines the splitting of the fourfold degeneracy of the subspace ($`Q=0`$, $`S=1/2`$, $`S_z=1/2`$) of the LM fixed point at energy $`ϵ_1+ϵ_3`$. We focus on the matrix element $`W_{12}`$: $$W_{12}=\psi _1\|H_N^{}\|\psi _2=\alpha (r)f(N)\psi _1\|\stackrel{}{S}_{\mathrm{imp}}\stackrel{}{s}_0\|\psi _2.$$ (54) The strategy for the calculations can be extended to the other matrix elements and the other subspaces, for which we add the perturbative results at the end of this appendix without derivation. The operator $`\stackrel{}{S}_{\mathrm{imp}}\stackrel{}{s}_0`$ is decomposed in four parts: $`\stackrel{}{S}_{\mathrm{imp}}\stackrel{}{s}_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}S_{\mathrm{imp}}^+c_0^{}c_0+{\displaystyle \frac{1}{2}}S_{\mathrm{imp}}^{}c_0^{}c_0`$ (55) $`+{\displaystyle \frac{1}{2}}S_{\mathrm{imp}}^z\left(c_0^{}c_0c_0^{}c_0\right),`$ so that $`W_{12}`$ can be written as $$W_{12}=\alpha (r)f(N)\frac{1}{2}\left[\mathrm{I}+\mathrm{II}+\mathrm{III}\mathrm{IV}\right],$$ (56) with $$\mathrm{I}=\psi _1\|S_{\mathrm{imp}}^+c_0^{}c_0\|\psi _2,$$ (57) and the other terms accordingly. With the definitions of $`\|\psi _1`$ and $`\|\psi _2`$ of eq. (LABEL:eq:four-psis) we have: $`\mathrm{I}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\psi _0\|\left(\xi _3^{}\xi _1+\xi _3^{}\xi _1\right)f_{}S_{\mathrm{imp}}^+c_0^{}c_0`$ $`\times \left[{\displaystyle \frac{1}{\sqrt{6}}}f_{}^{}\left(\xi _1^{}\xi _3\xi _1^{}\xi _3\right)+{\displaystyle \frac{2}{\sqrt{6}}}f_{}^{}\xi _1^{}\xi _3\right]\|\psi _0.`$ With $`S_{\mathrm{imp}}^+=f_{}^{}f_{}`$ we immediately see that the terms containing $`f_{}S_{\mathrm{imp}}^+f_{}^{}`$ drop out. The remaining impurity operators, $`f_{}S_{\mathrm{imp}}^+f_{}^{}`$, give unity when acting on $`\|\psi _0`$ so one arrives at $$\mathrm{I}=\frac{1}{\sqrt{3}}\left[\mathrm{Ia}+\mathrm{Ib}\right],$$ (59) with $`\mathrm{Ia}`$ $`=`$ $`\psi _0\|\xi _3^{}\xi _1c_0^{}c_0\xi _1^{}\xi _3\|\psi _0`$ $`\mathrm{Ib}`$ $`=`$ $`\psi _0\|\xi _3^{}\xi _1c_0^{}c_0\xi _1^{}\xi _3\|\psi _0.`$ (60) To analyze Ia and Ib, the operators $`c_{0\sigma }^{()}`$ have to be expressed in terms of the operators $`\xi _{p\sigma }^{()}`$: $$c_{0\sigma }=\underset{p^{}}{}\alpha _{0p^{}}\xi _{p^{}\sigma },c_{0\sigma }^{}=\underset{p}{}\alpha _{0p}^{}\xi _{p\sigma }^{},$$ (61) with the sums over $`p`$ and $`p^{}`$ defined in eq. (20). This gives $$\mathrm{Ia}=\underset{pp^{}}{}\alpha _{0p}^{}\alpha _{0p^{}}\psi _0\|\xi _3^{}\xi _1\xi _p^{}\xi _p^{}\xi _1^{}\xi _3\|\psi _0.$$ (62) The only non-zero matrix elements of eq. (62) are for $`p=p^{}=3`$: $`\mathrm{Ia}`$ $`=`$ $`\alpha _{03}^{}\alpha _{03}\psi _0\|\xi _3^{}\xi _1\xi _3^{}\xi _3\xi _1^{}\xi _3\|\psi _0`$ (63) $`=`$ $`\|\alpha _{03}\|^2.`$ Similarly, the term Ib gives $`\mathrm{Ib}`$ $`=`$ $`{\displaystyle \underset{pp^{}}{}}\alpha _{0p}^{}\alpha _{0p^{}}\psi _0\|\xi _3^{}\xi _1\xi _p^{}\xi _p^{}\xi _1^{}\xi _3\|\psi _0`$ (64) $`=`$ $`\|\alpha _{01}\|^2,`$ so that in total: $$\mathrm{I}=\frac{1}{\sqrt{3}}\left[\|\alpha _{03}\|^2+\|\alpha _{01}\|^2\right].$$ (65) The next term $`\mathrm{II}=\psi _1\|S_{\mathrm{imp}}^{}c_0^{}c_0\|\psi _2`$ gives zero due to the combination of impurity operators: $`f_{}f_{}^{}f_{}\mathrm{}`$ with $`f_{}`$ from $`\psi _1\|`$ and $`f_{}^{}f_{}=S_{\mathrm{imp}}^{}`$. The third term $`\mathrm{III}=\psi _1\|S_{\mathrm{imp}}^zc_0^{}c_0\|\psi _2`$ gives $`\mathrm{III}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{12}}}\psi _0\|\left(\xi _3^{}\xi _1+\xi _3^{}\xi _1\right)f_{}S_{\mathrm{imp}}^zc_0^{}c_0f_{}^{}`$ (66) $`\times \left(\xi _1^{}\xi _3\xi _1^{}\xi _3\right)\|\psi _0,`$ where the term with $`\frac{2}{\sqrt{6}}f_{}^{}\xi _1^{}\xi _3`$ from $`\|\psi _2`$ has already been dropped. So we are left with four terms $$\mathrm{III}=\frac{1}{\sqrt{12}}\left[\mathrm{IIIa}\mathrm{IIIb}+\mathrm{IIIc}\mathrm{IIId}\right],$$ (67) with $`\mathrm{IIIa}`$ $`=`$ $`\psi _0\|\xi _3^{}\xi _1f_{}S_{\mathrm{imp}}^zc_0^{}c_0f_{}^{}\xi _1^{}\xi _3\|\psi _0,`$ $`\mathrm{IIIb}`$ $`=`$ $`\psi _0\|\xi _3^{}\xi _1f_{}S_{\mathrm{imp}}^zc_0^{}c_0f_{}^{}\xi _1^{}\xi _3\|\psi _0,`$ $`\mathrm{IIIc}`$ $`=`$ $`\psi _0\|\xi _3^{}\xi _1f_{}S_{\mathrm{imp}}^zc_0^{}c_0f_{}^{}\xi _1^{}\xi _3\|\psi _0,`$ $`\mathrm{IIId}`$ $`=`$ $`\psi _0\|\xi _3^{}\xi _1f_{}S_{\mathrm{imp}}^zc_0^{}c_0f_{}^{}\xi _1^{}\xi _3\|\psi _0.`$ Following similar arguments as above one obtains $$\mathrm{IIIa}=\frac{1}{2}\underset{p}{}^{}\|\alpha _{0p}\|^2,$$ (69) where the $`p`$ in $`_{p}^{}{}_{}{}^{}`$ takes the values $$p=1,1,5,7,\mathrm{}N,$$ then $$\mathrm{IIIb}=\mathrm{IIIc}=0,$$ (70) and $$\mathrm{IIId}=\frac{1}{2}\underset{p}{}^{\prime \prime }\|\alpha _{0p}\|^2,$$ (71) where the $`p`$ in $`_{p}^{}{}_{}{}^{\prime \prime }`$ takes the values $$p=1,3,5,7,\mathrm{}N.$$ This gives for the third term $`\mathrm{III}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{12}}}\left[\mathrm{IIIa}\mathrm{IIId}\right]`$ (72) $`=`$ $`{\displaystyle \frac{1}{4\sqrt{3}}}\left[{\displaystyle \underset{p}{}^{}}\|\alpha _{0p}\|^2{\displaystyle \underset{p}{}^{\prime \prime }}\|\alpha _{0p}\|^2\right]`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{3}}}\left[\|\alpha _{01}\|^2\|\alpha _{03}\|^2\right].`$ The calculation of IV proceeds very similarly to III and one obtains $$\mathrm{III}=\mathrm{IV},$$ (73) so that we finally arrive at $`W_{12}`$ $`=`$ $`\alpha (r)f(N){\displaystyle \frac{1}{2}}\left(\|\alpha _{01}\|^2\|\alpha _{03}\|^2\right)\left[{\displaystyle \frac{1}{\sqrt{3}}}+0+2{\displaystyle \frac{1}{4\sqrt{3}}}\right]`$ (74) $`=`$ $`\alpha (r)f(N){\displaystyle \frac{1}{4}}\sqrt{3}\left(\|\alpha _{01}\|^2\|\alpha _{03}\|^2\right).`$ We performed a similar analysis for a couple of other subspaces. Here we list the results from the perturbative analysis for three more subspaces together with the corresponding basis states. ### A.1 $`Q=0`$, $`S=1/2`$, $`S_z=1/2`$, $`E=2ϵ_1`$ This subspace has the same quantum numbers $`Q`$, $`S`$ and $`S_z`$ as the one discussed above, so that the details of the calculation are very similar. The differences originate from the position of particles and holes in the single-particle spectrum of Fig. 6. This reduces the dimensionality of the subspace from four to two. The corresponding basis can be written as $`\|\psi _1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f_{}^{}(\xi _1^{}\xi _1+\xi _1^{}\xi _1)\|\psi _0,`$ $`\|\psi _2`$ $`=`$ $`\left[{\displaystyle \frac{1}{\sqrt{6}}}f_{}^{}(\xi _1^{}\xi _1\xi _1^{}\xi _1)+{\displaystyle \frac{2}{\sqrt{6}}}f_{}^{}\xi _1^{}\xi _1\right]\|\psi _0.`$ The first-order corrections are given by the 2$`\times `$2 matrix $$\{W_{ij}\}=\alpha (r)f(N)\left[\begin{array}{cc}0& \frac{\sqrt{3}}{4}\gamma \\ \frac{\sqrt{3}}{4}\gamma & \frac{1}{2}\beta \end{array}\right],$$ (76) with $`\gamma =\|\alpha _{01}\|^2\|\alpha _{01}\|^2`$ and $`\beta =\|\alpha _{01}\|^2+\|\alpha _{01}\|^2`$. Due to the particle-hole symmetry of the conduction band we have $`\|\alpha _{01}\|=\|\alpha _{01}\|`$; therefore, the off-diagonal matrix elements vanish and the effect of the perturbation is simply a negative energy-shift only for the state $`\|\psi _2`$: $$\{W_{ij}\}=\alpha (r)f(N)\left[\begin{array}{cc}0& 0\\ 0& \|\alpha _{01}\|^2\end{array}\right].$$ (77) This effect can be seen in the energy splitting of the first two low-lying excitations in Fig. 8. ### A.2 $`Q=1`$, $`S=0`$, $`E=ϵ_1`$ There is only one configuration for this combination of quantum numbers and excitation energy: $$\|\psi =\frac{1}{\sqrt{2}}(f_{}^{}\xi _1+f_{}^{}\xi _1)\|\psi _0.$$ (78) The first-order perturbation keeps the state in this one-dimensional subspace and the energy correction is given by $$\mathrm{\Delta }E=\psi \|H_N^{}\|\psi =\frac{3}{4}\alpha (r)f(N)\|\alpha _{01}\|^2.$$ (79) ### A.3 $`Q=1`$, $`S=0`$, $`E=ϵ_3`$ The difference to the previous case is the position of the hole in the single-particle spectrum. The state is now given by $$\|\psi =\frac{1}{\sqrt{2}}(f_{}^{}\xi _3+f_{}^{}\xi _3)\|\psi _0,$$ (80) with the energy correction $$\mathrm{\Delta }E=\psi \|H_N^{}\|\psi =\frac{3}{4}\alpha (r)f(N)\|\alpha _{03}\|^2.$$ (81) ## Appendix B Details of the Perturbative Analysis around the Strong Coupling Fixed Point The main difference in the calculation of the matrix elements $`\{W_{ij}\}`$ for this case is due to the structure of the perturbation, see eq. (41). Furthermore the ground state of the SC fixed point is fourfold degenerate and the perturbation partially splits this degeneracy, as discussed in the following. ### B.1 $`Q=0`$, $`S=1/2`$, $`S_z=1/2`$, $`E=0`$ This is one of the four degenerate ground states at the SC fixed point: $$\|\psi _1=\xi _0^{}\|\psi _0,$$ (82) with $`\|\psi _0`$ defined in eq. (45). The perturbative correction is given by $$\psi _1\|H_N^{}\|\psi _1=\frac{1}{2}\beta (r)\overline{f}(N)(1\|\alpha _{f0}\|^4)$$ (83) which corresponds to the energy shift of the ground state: $$\mathrm{\Delta }E_1=\frac{1}{2}\beta (r)\overline{f}(N)(1\|\alpha _{f0}\|^4).$$ (84) The coefficients $`\alpha _{fl}`$ are defined via the relation between the operators $`f_\sigma ^{()}`$ and $`\xi _{l\sigma }^{()}`$: $$f_\sigma =\underset{l^{}}{}\alpha _{fl^{}}\xi _{l^{}\sigma },f_\sigma ^{}=\underset{l}{}\alpha _{fl}^{}\xi _{l\sigma }^{}.$$ (85) ### B.2 $`Q=1`$, $`S=0`$, $`E=0`$ This state is also a ground state in the $`U=0`$ case: $$\|\psi _2=\|\psi _0.$$ (86) The calculation of the first-order correction for $`\|\psi _2`$ gives $$\psi _2\|H_N^{}\|\psi _2=\frac{1}{2}\beta (r)\overline{f}(N)(1+\|\alpha _{f0}\|^4).$$ (87) This means that the ground state including the effect of the perturbation is given by $`\|\psi _1`$ in eq. (82) and the state $`\|\psi _2`$ appears as an excited state. For a comparison with the energy levels shown in the NRG flow diagrams, where the ground state energy is set to zero in each iteration, we subtract the perturbative correction of the ground state ($`\mathrm{\Delta }E_1`$) from the energies of all other excited states. Subtracting this energy shift from eq. (87) gives the net energy correction for the $`\|\psi _2`$ state: $$\mathrm{\Delta }E_2=\beta (r)\overline{f}(N)\|\alpha _{f0}\|^4.$$ (88) ### B.3 $`Q=1`$, $`S=0`$, $`E=ϵ_2`$ The state corresponding to this subspace is given by: $$\|\psi _3=\frac{1}{\sqrt{2}}(\xi _0^{}\xi _2+\xi _0^{}\xi _2)\|\psi _0.$$ (89) The first-order correction reads $$\psi _3\|H_N^{}\|\psi _3=\beta (r)\overline{f}(N)\left[\frac{1}{2}(1\|\alpha _{f0}\|^4)+3\|\alpha _{f0}\|^2\|\alpha _{f2}\|^2\right].$$ (90) Subtracting the energy correction for the ground state results in $$\mathrm{\Delta }E_3=3\beta (r)\overline{f}(N)\|\alpha _{f0}\|^2\|\alpha _{f2}\|^2.$$ (91) ### B.4 $`Q=1`$, $`S=0`$, $`E=ϵ_4`$ Similarly for the state $$\|\psi _4=\frac{1}{\sqrt{2}}(\xi _0^{}\xi _4+\xi _0^{}\xi _4)\|\psi _0,$$ (92) the first-order correction is given by $$\psi _4\|H_N^{}\|\psi _4=\beta (r)\overline{f}(N)\left[\frac{1}{2}(1\|\alpha _{f0}\|^4)+3\|\alpha _{f0}\|^2\|\alpha _{f4}\|^2\right],$$ and subtracting the energy correction for the ground state results in: $$\mathrm{\Delta }E_4=3\beta (r)\overline{f}(N)\|\alpha _{f0}\|^2\|\alpha _{f4}\|^2.$$ (93)
warning/0506/cond-mat0506797.html	ar5iv	text	# Sequence heterogeneity and the dynamics of molecular motors ## 1 Introduction The study of molecular motors has been transformed in recent years with the increasing use of single molecule experiments . In one key experiment an external force is applied to a molecular motor opposing its motion . Typically, as the force is increased, the velocity of the motor decreases until it is completely stalled. The behavior of the velocity as a function of force provides much information on the chemical cycle underlying the motion of the motor. For example, the stall force is a direct estimate of the force exerted by the molecular motor. The experiments have also motivated much theoretical work on the dynamics of the motors . Fitting the experimentally obtained velocity-force curves allows extraction of detailed information on the chemical cycle of the motor . Most theoretical studies have focused on motors which move on featureless, or periodic linear tracks . Such a description would be appropriate, for example, for kinesin which moves along a microtubule filament, which is periodic, using only ATP for its motion . However, in many cases the assumption of a periodic medium fails. Examples of motors which move on heterogeneous tracks include RNA polymerase which moves along DNA, ribosomes which move along mRNA, helicases which unwind DNA, exonucleases which turn double-stranded DNA into single-stranded DNA and many others. All these motors move along tracks which are inherently “disordered” or heterogeneous due to the underlying sequence of the linear template. Theoretically, molecular motors moving along disordered tracks have received much less attention . Another form of heterogeneity which has largely been ignored arises from the different chemical fuels which may be used by molecular motors to move along the track. For example, RNA polymerase uses different nucleotide triphosphates (NTP’s) which build the mRNA it produces, each supplying a different amount of chemical energy, to move along a DNA strand. A different “annealed” form of disorder (in contrast to the “quenched” disorder embodied in a particular nucleotide sequence) can be present even in molecular motors moving along perfectly periodic tracks in a solution containing several distinct types of chemical fuels. For example, it is known that kinesin can move using other nucleotide triphosphates (such as GTP) instead of ATP, albeit less efficiently . Recently, we have introduced a simplified model for molecular motors which allows the effects of disorder to be studied in considerable detail . We have focused so far on heterogeneous tracks and argued that near the stall force the dynamics of the motor is strongly affected by the heterogeneity embodied in a particular DNA or RNA sequence. Due to the “sequence disorder” on which the motor is moving (many DNA sequences have only short range correlations ) the displacement of the motor as a function of time ceases to be linear in time close enough to the stall force. The displacement becomes sublinear in time, growing as $`t^\mu `$, with $`\mu `$ varying continuously from $`1`$ to $`0`$ as the stall force is approached. As discussed below there are also anomalies in the diffusive spreading about the average motor position which extend even further below the stall force. In this paper we review some of these results, stressing several experiments which could be performed to test the predictions of the model. We also explore the effect of heterogeneous fuels on the motion of molecular motors. In it was suggested that inhomogeneous fuel concentrations could enhance significantly the regime near the stall force over which anomalous dynamics is observed. Here we study this type of disorder numerically and illustrate the dramatic effect of varying the concentrations of the different fuels used to power the motor. For motors moving along heterogeneous tracks we also discuss the dynamics in the extreme limit where detailed balance is violated and motors never take backward steps. Finally, we consider motors moving along a periodic substrate powered by different kinds of fuels. We discuss, for simple cases, the expected velocity of the motor as the relative proportion of two different types of fuel in the solution is varied. The behavior of more complicated models is also discussed. ## 2 The Model In this section we define the model used throughout the paper. We start with a special case of the general class of $`n`$-state models explored by Kolomeisky and Fisher and consider a “minimal” motor with only two internal states. This simplified model reproduces important features of previously studied systems and allows us to explore generic behavior in new situations in a minimal form. The model is easily generalized to account for heterogeneous fuels and tracks. When appropriate, we will mention how results are modified for general $`n`$-state models. The model has been introduced and studied in detail in and here we only review its basic properties. We begin by assuming a perfectly periodic substrate. The location along the one-dimensional track, $`x`$, is assumed to take a discrete set of values $`x_m`$, where $`m=0,1,2\mathrm{}`$ labels distinct $`a`$ and $`b`$ sites. Although not essential, we assume for simplicity that the distances between $`x_{m+1}x_m`$ and $`x_{m+2}x_{m+1}`$ are equal and set $`x_{m+2}x_m=2a_0`$, which is the size of a step taken by the motor after completing a chemical cycle such as hydrolysis of ATP. In general, as discussed in , the distance traveled by the motor between internal states may be different for different internal transitions. However, we do not expect such modifications to affect the long time behavior over a range of parameters near the stall force. The dynamics embodied in the model is shown schematically in Fig. 1. Internal states labeled by $`a`$ have an energy $`\epsilon =0`$ while internal states labeled by $`b`$ have a higher energy $`\epsilon =\mathrm{\Delta }\epsilon `$. The local detailed balance condition (in temperature units such that $`k_B=1`$) is satisfied by our choice of rate constants, $`w_a^{}`$ $`=`$ $`(\alpha e^{\mathrm{\Delta }\mu /T}+\omega )e^{\mathrm{\Delta }\epsilon /Tf/2T}`$ $`w_b^{}`$ $`=`$ $`(\alpha +\omega )e^{f/2T}`$ $`w_a^{}`$ $`=`$ $`(\alpha ^{}e^{\mathrm{\Delta }\mu /T}+\omega ^{})e^{\mathrm{\Delta }\epsilon /T+f/2T}`$ (1) $`w_b^{}`$ $`=`$ $`(\alpha ^{}+\omega ^{})e^{f/2T}.`$ Following Ref. , there are two parallel channels for the motion. The first, represented by contributions containing $`\alpha `$ and $`\alpha ^{}`$, arise from utilization of chemical energy biased by a chemical potential difference $`\mathrm{\Delta }\mu `$ between, say ATP and the products of hydrolysis ADP and P<sub>i</sub>. The second channel, represented by the terms containing $`\omega `$ and $`\omega ^{}`$, correspond to thermal transitions unassisted by the chemical energy. We assume that the externally applied force $`F`$ biases the motion in a particularly simple way (consistent with detailed balance) and define $`f=Fa_0`$. If the substrate lacks inversion symmetry (a necessary condition for directed motion, driven by $`\mathrm{\Delta }\mu `$, when $`f=0`$ ), we have $`\alpha ^{}\alpha `$ and $`\omega ^{}\omega `$. If the fuel is ATP, the chemical potential difference which drives the motion is $$\mathrm{\Delta }\mu =T\left[\mathrm{ln}\left(\frac{[ATP]}{[ADP][P_i]}\right)\mathrm{ln}\left(\frac{[ATP]_{\mathrm{eq}}}{[ADP]_{\mathrm{eq}}[P_i]_{\mathrm{eq}}}\right)\right],$$ (2) where the square brackets $`[\mathrm{}]`$ denote concentrations under experimental conditions and the brackets $`[\mathrm{}]_{\mathrm{eq}}`$ denote the corresponding concentrations at equilibrium. The rate constants in Eq. 1 define a set of differential equations for the probability $`P_n(t)`$ of being at site $`n`$ at time $`t`$. For odd $`n`$ one has $$\frac{dP_n(t)}{dt}=w_a^{}P_{n1}(t)+w_a^{}P_{n+1}(t)(w_b^{}+w_b^{})P_n(t),$$ (3) while for even $`n`$ $$\frac{dP_n(t)}{dt}=w_b^{}P_{n1}(t)+w_b^{}P_{n+1}(t)(w_a^{}+w_a^{})P_n(t).$$ (4) It is illuminating, especially when we consider rate constants which depend on the position along a heterogenous track, to study two limits of these equations. In the first the chemical potential difference $`\mathrm{\Delta }\mu `$ and the applied force $`f`$ are small compared to the energy difference $`\mathrm{\Delta }\epsilon `$ so that $`b`$ states relax quickly compared to $`a`$ states. This condition implies that $`(w_b^{}+w_b^{})(w_a^{}+w_a^{})`$ so that in the long time-limit to a good approximation the left hand side of Eq. 3 may be set to zero. Upon solving for $`P_n(t)`$ with $`n`$ odd and substituting into Eq. 4, we obtain differential equations just for the even sites $`{\displaystyle \frac{dP_n(t)}{dt}}`$ $`=`$ $`{\displaystyle \frac{w_b^{}w_a^{}P_{n+2}(t)+w_b^{}w_a^{}P_{n2}(t)(w_b^{}w_a^{}+w_b^{}w_a^{})P_n(t)}{(w_b^{}+w_b^{})}},`$ (5) $`(\mathrm{n}\mathrm{even}).`$ Similarly in the limit $`\mathrm{\Delta }\mu \mathrm{\Delta }\epsilon `$ (with $`f`$ near the stall force) the motor spends most if its time in $`b`$ states. Now, in the long time-limit to a good approximation the left hand side of Eq. 4 may be set to zero. The remaining differential equations for the odd sites read $`{\displaystyle \frac{dP_n(t)}{dt}}`$ $`=`$ $`{\displaystyle \frac{w_b^{}w_a^{}P_{n+2}(t)+w_b^{}w_a^{}P_{n2}(t)(w_b^{}w_a^{}+w_b^{}w_a^{})P_n(t)}{(w_a^{}+w_a^{})}},`$ (6) $`(\mathrm{n}\mathrm{odd}).`$ Note that in both limits the dynamics of the motors on long-times can be described by a random walker moving on an effective energy landscape associated with what is in general a non-equilibrium dynamics. Upon absorbing the denominator factors $`(w_a^{}+w_a^{})`$ and $`(w_b^{}+w_b^{})`$ into a rescaling of the rate constants in the numerator, the effective energy landscape can be read off from Eq. 5 or Eq. 6. One finds that the effective energy difference between two sites which are two monomers apart is given by $$E_{n+2}E_n\mathrm{\Delta }E=T\mathrm{ln}\left(\frac{w_a^{}w_b^{}}{w_a^{}w_b^{}}\right).$$ (7) For a periodic track, this leads to a tilted energy landscape (with tilt controlled by $`\mathrm{\Delta }\mu `$ and $`f`$) and an effective energy difference between $`2m`$-adjacent monomers $`2m\mathrm{\Delta }E`$. The tilted energy landscape leads to diffusion with drift on long time-scales and large length-scales. The effective energy landscape can also be obtained by assuming detailed balance and equating the rate asymmetry between two neighboring even sites to an effective energy difference $`\mathrm{\Delta }E`$. It is straightforward to verify from Eq. 7 with Eq. 1 that for a periodic substrate no net motion is generated when the external force $`f=0`$ and the chemical potential difference $`\mathrm{\Delta }\mu =0`$. Also, when there is directional symmetry in the transition rates, $`\alpha =\alpha ^{}`$, $`\omega =\omega ^{}`$, and $`f=0`$ no net motion is generated even when $`\mathrm{\Delta }\mu 0`$. Absent this symmetry, chemical energy can be converted to motion. These conditions are equivalent to those presented in for continuum models and are exhibited here in a minimal model. The effect of the externally applied force is simply to bias the motion of the motor in the direction in which it is applied. The velocity for a motor moving along a periodic track in the two limits discussed above can be obtained by taking the continuum limit and yields $`v=(w_a^{}w_b^{}w_a^{}w_b^{})/(w_b^{}+w_b^{})`$ for the limit specified by Eq. 5 and $`v=(w_a^{}w_b^{}w_a^{}w_b^{})/(w_a^{}+w_a^{})`$ for the limit specified by Eq. 6. More generally, for periodic rates, it is straightforward to calcualte the velocity, for example using Bloch eigenfunctions $`\|ke^{ikx}`$ and expanding the eigenvalues in the wavevector $`k`$. The linear term gives the velocity and the quadratic part the effective diffusion constant. For the velocity one finds $$v=\frac{w_a^{}w_b^{}w_a^{}w_b^{}}{w_a^{}+w_b^{}+w_a^{}+w_b^{}}.$$ (8) An alternative to studying the solution of the equations, which will be very useful throughout this paper, is to use Monte-Carlo simulations. The rates specified in Eq. (1) can be simulated using the following procedure: To make the simulation efficient we first normalize the entering or leaving rates for a site so that the largest one is unity. Then, at each step we choose with equal probability attempting to move the motor to the right or left on the lattice. Following this choice a random number is drawn from a uniform distribution in the interval $`[0,1]`$. The motor is moved in the chosen direction provided the random number is smaller than the corresponding rate. Thus, if a motor finds itself on site $`a`$ in Fig. 1 and $`w_a^{}>w_a^{},w_b^{},w_b^{}`$ is the largest rate it will (after the rescaling) move one step to the right with probability $`1/2`$, one step to the left with probability $`1/2\left(w_a^{}/w_a^{}\right)`$ and it will stay put with probability $`1/2[1(w_a^{}/w_a^{})]`$. Note that the probability of actually moving one step (right or left) to the $`b`$-sublattice during a particular attempt is $`1/2[1+(w_a^{}/w_a^{})]`$. Once the $`b`$-sublattice is reached the procedure is repeated with the rates $`w_b^{}`$ and $`w_b^{}`$ (note that the probabilities are still obtained by dividing by $`w_a^{}`$). To compare, for example, velocities for different choices of rates the overall number of attempts is rescaled at the end by the fastest rate (taken to be $`w_a^{}`$ in the above example). The same procedure is followed for both homogeneous and heterogeneous tracks. For the latter the largest rate is chosen from all possible hopping rates along the track. This protocol ensures relaxation to equilibrium in the absence of chemical or mechanics driving forces . In we have analyzed in detail the motion of the model when the track is not periodic. Such tracks arise naturally, for example, for motors such as RNA polymerase or helicases which move on DNA which has a well defined sequence. In this case the energy difference $`\mathrm{\Delta }E(m)`$ now becomes an explicit function of the location $`m`$ along the track, due to the dependance of the rates on the location on the track. To understand the energy which arises for heterogeneous tracks consider the “integrating out” procedure applied to the three sites shown in Fig. 2, where three distinct motor binding energies, $`E_1`$, $`E_2`$ and $`E_3`$, are indicated explicitly. We work in the limit $`\mathrm{\Delta }\mu \mathrm{\Delta }E`$, with $`\mathrm{\Delta }E=E_1E_2`$ or $`\mathrm{\Delta }E=E_3E_2`$, and $`f`$ close to the stall force, so that the approximation leading to Eq. 6 (“integrating out” site $`2`$) is appropriate. As rates for the heterogeneous cluster shown in Fig. 2, we take $`w_a^{}(13)`$ $`=`$ $`[\alpha (13)e^{\mathrm{\Delta }\mu (13)/T}+\omega (13)]e^{(E_2E_1)/Tf/2T}`$ $`w_b^{}(13)`$ $`=`$ $`[\alpha (13)+\omega (13)]e^{f/2T}`$ $`w_a^{}(13)`$ $`=`$ $`[\alpha ^{}(13)e^{\mathrm{\Delta }\mu (13)/T}+\omega ^{}(13)]e^{(E_2E_3)/T+f/2T}`$ (9) $`w_b^{}(13)`$ $`=`$ $`[\alpha ^{}(13)+\omega ^{}(13)]e^{f/2T}.`$ where the arguments “$`(13)`$” appended to the $`w`$’s, $`\alpha `$’s,$`\alpha ^{}`$’s, $`\omega `$’s and $`\omega ^{}`$’s simply mean that these are the heterogenous rates appropriate to the cluster $`123`$. These rates obey detailed balance conditions for the two channels, and have a similar dependence on $`\mathrm{\Delta }\mu (13)`$ and $`f`$ and various energy differences as the rates in Eq. 1. The notation $`\mathrm{\Delta }\mu (13)`$ indicates that the chemical potential difference could depend on which NTP (in the case of RNA polymerase) provides the energy for that particular step. Upon assuming fast relaxation of site $`2`$ in Fig. 2, from a formula similar to Eq. 7, $`\mathrm{\Delta }E_{13}`$ $`=`$ $`E_3E_1+2fT\mathrm{ln}\left[{\displaystyle \frac{(\alpha (13)e^{\mathrm{\Delta }\mu /T}+\omega (13))(\alpha ^{}(13)+\omega ^{}(13))}{(\alpha ^{}(13)e^{\mathrm{\Delta }\mu /T}+\omega ^{}(13))(\alpha (13)+\omega (13))}}\right],`$ (10) $``$ $`E_3E_1+2f+\eta _{13}`$ Eq. 10 illustrates the following important points, applicable to motor molecules on heterogenous tracks more generally: ($`a`$) if $`\mathrm{\Delta }\mu (13)=0`$, then $`E_{13}=E_3E_1+2f`$, with a similar formula for all neighboring pairs of odd sites. Thus, in the absence of chemical energy, we have a “random energy landscape” with bounded energy fluctuations; ($`b`$) if $`\alpha (13)=\alpha ^{}(13)`$ and $`\omega (13)=\omega ^{}(13)`$ (inversion symmetry), chemical energy does not lead to net motion between sites $`1`$ and $`3`$, as discussed above; ($`c`$) in general, $`\mathrm{\Delta }E_{13}=E_3E_1+2f+\eta _{13}`$, where $`\eta _{13}`$ is a random function of position along the heterogeneous track. Upon passing to a coarse-grained position $`\eta (m)`$, where $`m`$ is the position along the track we see that the effective ‘coarse grained’ energy difference between two points $`m_1`$ and $`m_2=m_1+m`$ which are $`m`$-monomers apart is given by $`_{m^{}=m_1}^{m_1+m}\eta (m^{})`$. The landscape itself behaves like a random walk with fluctuations which grow as $`\sqrt{m}`$, corresponding to a random-forcing energy landscape (for RNA polymerase, the different chemical potentials of the nucleotides in the transcript also contribute to a random forcing landscape ). The self-similar structure of the random force landscape leads to interesting dynamics near the stall force. As the stall force is approached the dynamics slows down and becomes dominated by motion between deep minima of the energy landscape . The minima correspond to specific locations along the track where the motor tends to pause. The distribution of dwell times at these minima, $`P(\tau )`$, averaged over the different locations on the track, is expected to behave as $`\tau ^{(1+\mu )}`$, where $`\mu `$ (not to be confused with a chemical potential!) is related to the force, fluctuations in the effective energy landscape and temperature. For random forcing energy landscape where the energy difference between two points is drawn from a Gaussian distribution with a variance $`V=\overline{\eta (m)^2}`$, where the overline denotes an average along the sequence, one can show that $$\mu (f)=2T\|\overline{\mathrm{\Delta }E}_{f=0}2f\|/V,$$ (11) where $`\overline{\mathrm{\Delta }E}_{f=0}`$ is the mean slope of the potential (averaged along the sequence) at zero force. The exponent $`\mu `$ thus decreases continuously to zero as $`f`$ increases toward the stall force of the motor (defined by $`\mu (f_s)0`$). For more general distributions of the effective energy difference the value of $`\mu `$ might be different from Eq. 11 by factors of order unity. Near the stall force the expected distribution of pause times becomes broader as $`\mu `$ becomes closer to $`0`$. The dynamics of the motors are altered from diffusion with drift when the pause-time distribution becomes very broad. The dynamics then depends on the numerical value of $`\mu `$ defined in Eq. 11 : * $`\mu <1`$ – Around the stall force, both the drift and diffusive behavior of the motor become anomalous. The displacement of the motor as a function of time increases as $`t^\mu `$. Thus, in this region the velocity is undefined, in the sense that it depends on the experimental observation time, $`t_E`$, through $`vt_E^{\mu 1}`$. Moreover, the spread of the probability distribution of the motor about its mean position also behaves anomalously with a variance which grows as $`t_E^{2\mu }`$. Experimentally, for a given $`t_E`$, this anomaly should lead to a convex velocity as a function of force curve in the vicinity of the stall force. The curve will become more and more convex as $`t_E`$ is increased; the velocity actually vanishes for a range of $`f`$’s near the stall force in the limit $`t_E\mathrm{}`$. * $`1<\mu <2`$ – Further away from the stall force the displacement of the motor as a function of time grows linearly. At long-times the velocity becomes independent of the averaging window. However, the variance of the probability distribution around the mean is anomalous and grows as $`t_E^{2/\mu }`$. * $`\mu >2`$ – Far below the stall force both the displacement and the variance of the probability distribution around the mean grow linearly in time, as in conventional diffusion with drift. These results can easily be shown to apply as well to general $`n`$-state models. Moreover, it can be argued that even if several parallel channels exist for moving from one monomer to another the results are also qualitatively unchanged . Experimentally, the predictions of the model can be tested by measuring the displacement of the motor as a function of time, averaged over different experimental runs (and, possibly, sequences). Each time trace of the motor position will have an irregular shape due to pauses, which will increase in duration as the stall force is approached, at specific locations along the track. However, averaged over many time traces (or sequences) the expected displacement will grows as $`t^\mu `$ with $`\mu <1`$ close enough to the stall force. Note that if one averages over time traces for a fixed sequence, the displacement is expected to grow as $`s(t)t^\mu `$ where $`s(t)`$ has fluctuations of order unity, because the sequence information is not completely erased in this case. An alternative experimental test would be to measure the distribution of dwell times $`𝒫(\tau )`$. Because $`𝒫(\tau )1/\tau ^{1+\mu }`$ for large $`\tau `$, the distribution becomes wider as the stall force is approached. Monitoring $`𝒫(\tau )`$ has the advantage of probing the wide distribution even in regimes which are not very close to the stall force. We now consider limiting cases of the above model on heterogeneous tracks. These illustrate a number of interesting features and suggest ways in which the anomalous dynamics might be observed experimentally. We will also use the model to explore heterogeneous chemical energy sources for motors on a periodic track. This situation may be realized in motors such as kinesin which can use several types of chemical energy to move along the track. From the two state model described above we deduce the expected behavior of the velocity as the relative proportions of the different fuels is varied and mention generalizations to more general $`n`$-state models. ## 3 Strongly biased motors In this section we study motors moving on a heterogeneous substrate in the limit where one of the transition rates is strongly biased in a certain direction. An extreme limit occurs when one of the transition rates, in, say, the backward direction, is zero. Although this limit violates detailed balance, it could be a reasonable approximation for certain strongly biased experiments. One such model is a special case of the two state model discussed in Section 2: $`w_a^{}`$ $`=`$ $`(\alpha e^{\mathrm{\Delta }\mu /T}+\omega )e^{\mathrm{\Delta }\epsilon /Tf/2T}`$ $`w_b^{}`$ $`=`$ $`(\alpha +\omega )e^{f/2T}`$ $`w_a^{}`$ $`=`$ $`\omega ^{}e^{\mathrm{\Delta }\epsilon /T+f/2T}`$ (12) $`w_b^{}`$ $`=`$ $`\omega ^{}e^{f/2T}.`$ Here we have set $`\alpha ^{}=0`$ so that in the $`\mathrm{\Delta }\mu T`$ limit the motor will be strongly biased to move towards the right. Physically, this situation corresponds to a motor which can use chemical energy only to move in a certain direction. (We expect qualitatively similar results for a wide variety of strongly biased “one way” models). Next, we assume an extremely strong bias $`\mathrm{\Delta }\mu T`$ limit of the model such that $`\alpha e^{\mathrm{\Delta }\mu /T}`$ is very large but $`\alpha `$ and $`\omega `$ are so small that $`w_b^{}`$ can be set to be zero. This limit only makes sense far from the stall force. The stall force of a model with $`w_b^{}=0`$ (as any other model were one of the reactions is assumed to be unidirectional) is infinite: The effective energy landscape, Eq. 7, which describes the dynamics of such a limit has an infinite slope. In this section we will compare numerically trajectories when $`w_b^{}=0`$ and when $`w_b^{}0`$. Using Eq. 11, the infinite slope of the effective energy landscape implies that $`\mu =\mathrm{}>2`$, so that the dynamics is diffusion with drift. The linear drift is illustrated clearly in Fig. 3 where trajectories of the strongly biased model are shown for different forces. Note that even for very large forces the displacement of the motor as a function of time grows linearly. It can be shown that the dwell time distribution on the track decays exponentially $`𝒫(\tau )e^{\tau /\tau ^{}}`$ in contrast to the power law distribution expected when $`w_b^{}0`$. To observe any significant pausing clearly one must have $`f\mathrm{\Delta }\mu `$ (see inset of Fig. 3). In this limit, however, $`w_a^{}w_a^{}`$, and backwards motion cannot be neglected (see Eq. 12). In Fig. 4 we show the very different behavior of simulations where we take the back hopping to be non-zero. Now we do not neglect $`\omega `$ and the backward hopping rate $`w_b^{}`$ is nonzero! For forces near the stall force the displacement of the motor as a function of time seems to saturate even for a trajectory generated by a single numerical experiment for a particular sequence. Such a behavior is consistent with that expected from a sublinear displacement of the motor. Note, however, that with the exception of the largest force, all curves on long enough time scales are expected to yield, after an average over many thermal realizations, an asymptotically linear curve (see the corresponding values of $`\mu `$ presented in the figure). However, even for these curves, pausing is pronounced despite the fact that for small resisting forces the motor rarely moves backwards. ## 4 Effect of heterogeneous fuel concentrations for a motor moving along a heterogeneous substrate For some molecular motors the type of fuel which is used to move depends on the specific site along the track. For example, in the case of RNA polymerase, which produces messenger RNA, the energy from the hydrolysis of the specific NTP which is added to the mRNA chain is used for motion. While a random forcing energy landscape would exist even if the chemical energy released from every NTP were the same (see Sec. 2), the different chemical energies enhance the variance of the slopes, $`V`$, of the random forcing energy landscape (due to the different chemical potentials of the nucleotides in the transcript). This in turn (see Eq. 11) lowers the value of $`\mu `$ as compared to the case of equal chemical energies. As suggested in the variance $`V`$ could be further increased by increasing the concentration difference between the different NTPs in the solution. In an extreme situation, where one of the NTPs is completely removed from the solution, the motor will stall at specific locations where the NTP is needed. This trick is used in experiments to synchronize and control the motion of RNA polymerases . In this section we illustrate using the simple model, Eq. 1, the effects of changing NTP concentrations on motor motion. Although we work within a “minimal model” we expect the same effects for more complicated models of molecular motors with many internal states (for an example of such a model for RNA polymerase see ). To this end, we studied motors which can use two kinds of fuels depending on its location along the track. We hold the concentration of one fuel fixed and lower the concentration of the other by reducing the chemical potential associated with it. For reference we also study the case when the fuels have equal chemical potentials. Fig. 5 displays results of numerical simulations of the model defined by Eq. (12). Similar results were obtained with the more general model of Eq. (1). We show simulations with $`f/T=0`$ and when for one fuel $`e^{\mathrm{\Delta }\mu /T}=500`$ and for the other $`e^{\mathrm{\Delta }\mu /T}=500,50`$ or $`5`$. From a generalization of Eq. 2, we see that this latter variation corresponds to a change of two orders of magnitude in the concentration of the second fuel . As can be seen from Fig. 5, the effect of reducing one of the chemical potentials is to slow down the motion. However, note that for all fuel concentrations the motor displacement as a function of time is linear with these parameter values. Since at zero applied force one expects the motor always to be biased preferentially in a given direction which is independent of the monomer this behavior will hold even for larger differences in concentration. The change in concentration of one of the fuels can make the regime of anomalous dynamics much larger. In Fig. 6 we show the result of applying a force which opposes the motion of the motors on the dynamics. As can be easily seen from the figure, the velocity of the motors is, of course, lower in comparison to that when no force is applied. However, note that the motion of the curve with the largest difference in chemical potential (barely visible at the bottom) is very different. The motor almost immediately stalls after it starts moving. This example illustrates clearly how changing the concentrations of the different fuels can make the regime of anomalous dynamics easy to access at lower forces. In fact, averaged over many thermal and sequence realizations, the curve with the largest difference in chemical potentials shows a displacement of the motor which grows sublinearly in time, indicating that $`\mu <1`$. ## 5 Heterogeneous fuels on periodic tracks In this section we consider a molecular motor which is moving along a periodic track in a solution which contains more than one type of molecule which can supply it with chemical energy. The situation arises for kinesin in the presence of both ATP and, say, GTP. It is known that, while less efficient, alternative NTP molecules can also be used by kinesin to move along the track . We consider, for simplicity, a solution with two types of chemical fuels within the simple two state model (The analysis presented can easily be generalized to include additional fuel types). In addition, we generalize our model to treat two separate cases. In the first we assume that the internal states of the motor (i.e., the $`a`$ and $`b`$ sites in Fig. 1) are independent of the fuel used so that the chemical fuels are only used to move between states with fuel-dependent potential differences driving the changes. This situation arises when the fuel from the motor is used and released so quickly that the motor is unbound to the fuel in the “excited” internal state. In this case the fuel binding to the motor does not define an internal state. In the second case we study we allow for additional internal states of the motor, depending on which type of chemical fuel is bound to it. This situation arises when internal “excited” states include a fuel bound to the motor. Since the motor bound to the two fuels defines two distinct internal states, a direct thermal transition between them is therefore not possible. While both cases discussed below involve parallel pathways for transitions across a monomer, they are distinct in the type of internal states. ### 5.1 Case I: Internal states of the motors are independent of the chemical fuel This case amounts to a straightforward generalization of the rates in Eq. 1 to allow for multiple fuels. As mentioned above, we assume the internal states of the motor are independent of the chemical energy and that chemical energy is only used to assist the transitions. With two different chemical fuels, the generalized rates entering Eqs. 3 and Eq. 4 are now $`w_a^{}`$ $`=`$ $`(\alpha _1e^{\mathrm{\Delta }\mu _1/T}+\alpha _2e^{\mathrm{\Delta }\mu _2/T}+\omega )e^{\mathrm{\Delta }\epsilon /Tf/2T}`$ $`w_b^{}`$ $`=`$ $`(\alpha _1+\alpha _2+\omega )e^{f/2T}`$ $`w_a^{}`$ $`=`$ $`(\alpha _1^{}e^{\mathrm{\Delta }\mu _1/T}+\alpha _2^{}e^{\mathrm{\Delta }\mu _2/T}+\omega ^{})e^{\mathrm{\Delta }\epsilon /T+f/2T}`$ (13) $`w_b^{}`$ $`=`$ $`(\alpha _1^{}+\alpha _2^{}+\omega ^{})e^{f/2T}.`$ There are now three parallel paths, two assisted by chemical energy, one thermal. The subscript $`1`$ or $`2`$ refers to whether fuel “$`1`$” or “$`2`$” is being used for the transition. Standard relations for the chemical potential difference imply that $`e^{\mathrm{\Delta }\mu _i/T}`$ grows linearly with the concentration of fuel “$`i`$” (see Eq. 2). Because the extra channel is present, when one of the chemical potential differences is set to zero the model reduces to the single motor model but with modified rates as compared to a motor in a solution where the fuel is completely absent. For example, if $`\mathrm{\Delta }\mu _2=0`$ the model can be written in terms of Eq. 1 with the identifications $`\omega \omega +\alpha _2`$ and $`\omega ^{}\omega ^{}+\alpha _2^{}`$. Note also that, in contrast to disordered tracks where the motion of the motor is described by a random walker moving on a random force landscape, here the energy landscape is periodic except for a well defined tilt given by Eq. 7. With the help of Eq. 8, it is straightforward to verify that for this model the velocity takes the form $$v=\frac{A+B[x_1]+B^{}[x_2]}{C+D[x_1]+D^{}[x_2]},$$ (14) where $`[x_1]`$ and $`[x_2]`$ are the concentrations of fuel $`1`$ and $`2`$ respectively and $`A,B,B^{},C,D,D^{}`$ are complicated functions of the the coefficients in Eq. 13. The velocity thus varies as ratio of two polynomials, each depending linearly on the concentration of both fuels. Consider, for example, the velocity in an experiment where one of the fuels is held at constant concentration and the concentration of the other is varied. The external force is set to $`f=0`$; (varying $`f`$ did not affect the qualitative features discussed below). As shown in Fig. 7, the average velocity smoothly crosses over from a small value to a larger value in a sigmoidal fashion as the concentration of more energy rich fuel (“fuel $`1`$”) is increased. Note that $`\mathrm{\Delta }\mu _1/T`$ varies logarithmically with the concentration of fuel number $`1`$ . Thus, the concentration of fuel “$`1`$” varies over many orders of magnitude for the range of $`\mathrm{\Delta }\mu _1/T`$ shown in Fig. 7. In typical experiments, only a small portion of this crossover may be visible. For small $`\mathrm{\Delta }\mu _1/T`$, motor movement is controlled by fuel “$`2`$” while for large $`\mathrm{\Delta }\mu _1/T`$ it is controlled by fuel “$`1`$”. An analogous plot for motor velocity vs. $`\mathrm{\Delta }\mu _1/T`$ when fuel $`2`$ is absent entirely is shown for reference in Fig. 8. As expected, the velocity no longer exhibits a sigmoidal crossover between two regimes. ### 5.2 Case II: Internal states of the motors are coupled to the chemical fuel Next, we consider a different model incorporating two fuels. We now assume that the type of fuel molecule used to move the motor determines the entire chemical cycle which leads the motor to move across one monomer. An example might be the $`n=2`$ motor landscape shown in Fig. 1 where one step (e.g.; $`a`$ site $``$ $`b`$ site) involves the fuel molecule binding to the motor. In this case the entire sequence of transitions would be dictated by the fuel that is used. Thus, the choice of rates (either in the forward or backward direction) would be dictated by the initial step which is chosen at random, depending on the type of fuel utilized. Since the state of the motor is directly coupled to the chemical fuel, a pure thermal transition between the states is not possible (unless it involves moving across the monomer through parallel pathways not related to the motor, an effect which is ignored here). The new feature is that the motor can move across a monomer by choosing one of two distinct chemical pathways. This is in contrast to the model defined by Eq. 13 where two distinct chemical pathways exist for passing between the internal states of the motor, but where the transition across the monomer can occur via a mixture different chemical (or thermal) pathways. An example for such a model, where each fuel is modeled by a distinct channel which is a special case of the two state model defined by (1) is given by: $`w_a^{}`$ $`=`$ $`\alpha e^{\mathrm{\Delta }\mu _1/Tf/2T}`$ $`w_b^{}`$ $`=`$ $`\alpha e^{f/2T}`$ $`w_a^{}`$ $`=`$ $`\omega e^{f/2T}`$ (15) $`w_b^{}`$ $`=`$ $`\omega e^{f/2T},`$ for the first channel and $`u_a^{}`$ $`=`$ $`\gamma e^{\mathrm{\Delta }\mu _2/Tf/2T}`$ $`u_b^{}`$ $`=`$ $`\gamma e^{f/2T}`$ $`u_a^{}`$ $`=`$ $`\nu e^{f/2T}`$ (16) $`u_b^{}`$ $`=`$ $`\nu e^{f/2T}.`$ for the second channel. The fuel concentrations enter through the chemical potential differences $`\mathrm{\Delta }\mu _1`$ and $`\mathrm{\Delta }\mu _2`$ (see Eq. 2). The relative fuel abundances therefore control the ratio of the rates $`w_a^{}`$ and $`u_a^{}`$. The model will be realized physically in a motor which binds a fuel and uses its chemical energy in the transitions $`w_a^{}`$ or $`u_a^{}`$. The transitions $`w_b^{}`$ and $`u_b^{}`$ involve the release of the corresponding fuel. Here for simplicity we have set the energy difference between the states to be zero and neglected parallel thermal channels. We do not expect the qualitative results described below to be affected by such complications. The velocity of the model can be calculated in a straightforward manner and is found to be $$v=\frac{(u_b^{}+u_b^{})(w_a^{}w_b^{}w_a^{}w_b^{})+(w_b^{}+w_b^{})(u_a^{}u_b^{}u_a^{}u_b^{})}{(u_b^{}+u_b^{})(w_b^{}+w_b^{})+(u_b^{}+u_b^{})(w_a^{}+w_a^{})+(u_a^{}+u_a^{})(w_b^{}+w_b^{})}.$$ (17) Note that as in Case I, the velocity behaves as a ratio of two polynomials which are linear in the $`e^{\mathrm{\Delta }\mu _i/T}`$, and hence with each of the fuel concentrations (similar to Eq. 14). Also similar is the fact that the presence of a second channel alters the velocity even when one of the chemical potential difference is zero. Again, due to the presence of the second channel the velocity is lower than that of the motor in the presence of a single fuel. We do not present plots of the resulting velocity as the concentration of one of the fuels is varied since the qualitative features are similar to those presented in the previous subsection. Finally, we comment that more complicated scenarios (for example, the existence of additional parallel pathways, or more general $`n`$-state models) might change the explicit dependence on the concentrations of the fuels. However, we expect a general form of a ratio of two polynomials in the concentrations of the fuels even in much more complicated scenarios. Indeed, with a specific experiment in mind and some structural information on the motor one might be able to use such experiments, coupled with an analysis similar to that presented above, to deduce the number of steps in the chemical cycle which depend explicitly on the fuel concentration. For example, we have considered models where the internal states are independent of the chemical fuel and found that the general velocity can be a ratio of polynomials of a degree which is related to the number of steps which depend on the chemical fuel. However, the detailed results were very dependent on the choice of rates made. Acknowledgments: We are very grateful to D. K. Lubensky for many stimulating conversations and the collaboration which led to references and . We also thank L. Bau, M. D. Wang and K. C. Neuman for useful conversations and J. Gelles and N. Guydosh for interesting us in the problem of different fuels for kinesin on periodic tracks like microtubules. D.R.N was supported by the National Science Foundation through Grant No. DMR-0231631 and the Harvard Materials Research Laboratory via Grant No. DMR-0213805. Y.K was supported by the Human Frontiers Science Program. Permanent address: Department of Physics, Technion, Haifa 32000, Israel.
warning/0506/math-ph0506077.html	ar5iv	text	# A first–order purely frame–formulation of General Relativity ## Abstract In the gauge natural bundle framework a new space is introduced and a first–order purely frame–formulation of General Relativity is obtained. PACS number: 04.20.Fy, 11.10.-z Mathematics Subject Classification: 70S99, 83C99 Keywords: gauge natural theories, General Relativity, tetrad, variational calculus In some of our recent works a new geometrical framework for Yang–Mills field theories and General Relativity in the tetrad–affine formulation has been developed. The construction of the new geometrical setting has been obtained quotienting the first–jet bundles of the configuration spaces of the above theories in a suitable way, resulting into the introduction of a new family of fiber bundles. In this letter we show that these new spaces allow a (covariant) first–order purely frame–formulation of General Relativity. The whole geometrical construction will be developed within the gauge natural bundle framework , which provides the suitable mathematical setting for globally describing gravity in the tetrad formalism. To start with, let $`M`$ be a space–time manifold, allowing a metric tensor $`g`$ with signature $`\eta =(1,3)`$: the manifold $`M`$ will be called a $`\eta `$-manifold and the metric tensor canonical representation will be $`\eta ^{\mu \nu }:=diag(1,1,1,1)`$. Moreover, let $`L(M)`$ be the frame–bundle over $`M`$ and $`PM`$ a principal fiber bundle over $`M`$ with structural group $`SO(1,3)`$. The configuration space of the theory (the tetrad space) is a $`GL(4,\mathrm{})`$ bundle $`\pi :M`$, associated to $`P\times _ML(M)`$ through the left–action $$\lambda :(SO(1,3)\times GL(4,\mathrm{}))\times GL(4,\mathrm{})GL(4,\mathrm{}),\lambda (\mathrm{\Lambda },J;X)=\mathrm{\Lambda }XJ^1$$ (1) Taking eq. (1) into account, the space $``$ can be referred to local fibered coordinates $`x^i,e_i^\mu `$ $`(i,\mu =1\mathrm{}4)`$, undergoing the transformations laws $$\overline{x}^j=\overline{x}^j(x^i),\overline{e}_j^\mu =e_i^\sigma \mathrm{\Lambda }_\sigma ^\mu (x)\frac{x^i}{\overline{x}^j}$$ (2) where $`\mathrm{\Lambda }_\sigma ^\mu (x)SO(1,3)xM`$. Under these circumstances the tetrad fields can be identified with the sections of the bundle $`M`$. It is worth noticing that the conditions making $`M`$ into a $`\eta `$-manifold allow to choose the principal bundle $`P`$ in such a way that $``$ admits global sections (see ). In the following, such a choice will be systematically adopted. Moreover, we also remind that (compare with again) there exists a one-to-one correspondence between the global sections of the bundle $`\pi :M`$ and the principal morphisms $`i:PL(M)`$. Whenever two principal connections $`\omega _{i\nu }^\mu `$ over $`P`$ and $`\mathrm{\Gamma }_{ih}^k`$ over $`L(M)`$ are given, the covariant exterior differential of any tetrad field $`e^\mu (x)=e_i^\mu (x)dx^i`$ is well defined as $$De^\mu :=_je_i^\mu dx^jdx^i$$ (3) where $$_je_i^\mu =\frac{e_i^\mu }{x^j}+\omega _{j\nu }^\mu e_i^\nu \mathrm{\Gamma }_{ji}^ke_k^\mu $$ The first jet bundle associated to the fibration $`\pi :M`$ is now taken into account. A set of jet–coordinates over $`j_1()`$ is provided by $`x^i,e_i^\mu ,e_{ij}^\mu \left(\frac{e_i^\mu }{x^j}\right)`$, subject to the transformation laws (2) together with $$\overline{e}_{jk}^\mu =e_{ih}^\sigma \frac{x^h}{\overline{x}^k}\mathrm{\Lambda }_\sigma ^\mu \frac{x^i}{\overline{x}^j}+e_i^\sigma \frac{\mathrm{\Lambda }_\sigma ^\mu }{x^h}\frac{x^h}{\overline{x}^k}\frac{x^i}{\overline{x}^j}+e_i^\sigma \mathrm{\Lambda }_\sigma ^\mu \frac{^2x^i}{\overline{x}^k\overline{x}^j}$$ (4) The frame–formulation of general relativity that we propose here is based on the introduction of the following equivalence relation on $`j_1()`$. Let $`z=(x^i,e_i^\mu ,e_{ij}^\mu )`$ and $`\widehat{z}=(x^i,\widehat{e}_i^\mu ,\widehat{e}_{ij}^\mu )`$ be two elements of $`j_1()`$, chosen in such a way that they have the same projection over $`M`$, namely $`\widehat{\pi }(z)=\widehat{\pi }(\widehat{z})=x`$, with $`\widehat{\pi }:j_1()M`$. We denote by $`e^\mu `$ and $`\widehat{e}^\mu `$ two different sections of the bundle $`\pi :M`$, respectively chosen among the representatives of the equivalence classes $`z`$ and $`\widehat{z}`$. Then, we make $`z`$ equivalent to $`\widehat{z}`$ if and only if $$e^\mu (x)=\widehat{e}^\mu (x)\mathrm{and}De^\mu (x)=D\widehat{e}^\mu (x)$$ (5) for every choice of a principal connection $`\omega `$ on $`P`$ and $`\mathrm{\Gamma }`$ on $`L(M)`$. It is easy to see that $`z\widehat{z}`$ if and only if the following local coordinates expression holds: $$e_i^\mu =\widehat{e}_i^\mu \mathrm{and}(e_{ij}^\mu e_{ji}^\mu )=(\widehat{e}_{ij}^\mu \widehat{e}_{ji}^\mu )$$ (6) We denote by $`𝒥()`$ the quotient space $`𝒥():=j_1(E)/`$ and by $`\rho :j_1(E)𝒥()`$ the corresponding quotient canonical projection. A system of local fibered coordinates on the bundle $`𝒥()`$ is provided by $`x^i,e_i^\mu ,E_{ij}^\mu :=\frac{1}{2}\left(e_{ij}^\mu e_{ji}^\mu \right)(i<j)`$, subject to the transformation laws (2), together with: $$\overline{E}_{jk}^\mu =E_{ih}^\sigma \mathrm{\Lambda }_\sigma ^\mu \frac{x^h}{\overline{x}^k}\frac{x^i}{\overline{x}^j}+\frac{1}{2}e_i^\sigma \frac{\mathrm{\Lambda }_\sigma ^\mu }{x^h}\frac{x^h}{\overline{x}^k}\frac{x^i}{\overline{x}^j}\frac{1}{2}e_i^\sigma \frac{\mathrm{\Lambda }_\sigma ^\mu }{x^h}\frac{x^h}{\overline{x}^j}\frac{x^i}{\overline{x}^k}$$ (7) The geometry of the quotient space $`𝒥()`$ has been deeply examined in some previous papers . As a matter of fact, the quotient projection endows the bundle $`𝒥()`$ of most of the standard features of jet–bundles geometry. The principal results are shortly reported below (see for a more detailed discussion). $``$ $`𝒥`$-extension of sections. The $`𝒥`$-extension of a section $`\sigma :M`$ is defined as $`𝒥\sigma :=\rho j_1\sigma `$, namely projecting the jet–extension $`j_1(\sigma )`$ on $`𝒥()`$ by means of the quotient projection $`\rho `$. A section $`\gamma :M𝒥()`$ is said holonomic if there exists a section $`\sigma :M`$ such that $`\gamma =𝒥\sigma `$. In local coordinates, a section $`\gamma `$ is holonomic if and only if $`\gamma :x\left(x^i,e_i^\mu (x),E_{ij}^\mu (x)=\frac{1}{2}\left(\frac{e_i^\mu (x)}{x^j}\frac{e_j^\mu (x)}{x^i}\right)\right)`$. $``$ Contact forms. Let us define the following $`2`$-form on $`𝒥()`$: $$\theta ^\mu :=de_i^\mu dx^i+E_{ij}^\mu dx^idx^j$$ (8) where $`E_{ij}^\mu :=E_{ji}^\mu `$ whenever $`i>j`$. Under a change of local coordinates (2) and (7), the $`2`$-forms (8) undergo the transformation laws $$\overline{\theta }^\mu =\mathrm{\Lambda }_\nu ^\mu \theta ^\nu $$ (9) The vector bundle which is locally spanned by the $`2`$-forms (8) will be called the contact bundle $`𝒞(𝒥())`$ and any section $`\eta :𝒥()𝒞(𝒥())`$ will be called a contact $`2`$-form. Contact forms $`\eta `$ are such that $`\gamma ^{}(\eta )=0`$ whenever $`\gamma :M𝒥()`$ is holonomic. Conversely, if a section $`\gamma :M𝒥()`$ is such that $`\gamma ^{}(\eta )=0`$ for all contact forms $`\eta `$, then $`\gamma `$ is holonomic. $``$ $`𝒥`$-prolongations of morphisms and vector fields. A suitable family of morphisms $`\mathrm{\Phi }:`$, fibered over $`M`$, can be raised to a family of morphisms $`𝒥\mathrm{\Phi }:𝒥()𝒥()`$ considering their ordinary jet–prolongations and projecting them to $`𝒥()`$ through the quotient map, namely: $$𝒥\mathrm{\Phi }(z):=\rho j_1\mathrm{\Phi }(w)w\rho ^1(z),z𝒥()$$ In order that the above definition makes sense, such morphisms $`\mathrm{\Phi }:`$ have to satisfy the condition: $$\rho j_1\mathrm{\Phi }(w_1)=\rho j_1\mathrm{\Phi }(w_2)w_1,w_2\rho ^1(z),z𝒥()$$ (10) Referring to for the proof, it is easy to see that the only morphisms satisfying condition (10) are necessarily of the form: $$\{\begin{array}{c}y^i=\chi ^i(x^j)\hfill \\ \\ \widehat{e}_i^\nu =\mathrm{\Phi }_i^\nu (x^j,e_j^\mu )=\mathrm{\Gamma }_\mu ^\nu (x)\frac{x^r}{y^i}e_r^\mu +f_i^\nu (x)\hfill \end{array}$$ (11) where $`\mathrm{\Gamma }_\mu ^\nu (x)`$ and $`f_i^\nu (x)`$ are arbitrary local functions on $`M`$. Their $`𝒥`$prolongation is: $$\{\begin{array}{c}y^i=\chi ^i(x^k)\hfill \\ \\ \widehat{e}_i^\nu =\mathrm{\Gamma }_\mu ^\nu (x)\frac{x^r}{y^i}e_r^\mu +f_i^\nu (x)\hfill \\ \\ \widehat{E}_{ij}^\nu =\mathrm{\Gamma }_\mu ^\nu E_{ks}^\mu \frac{x^k}{y^i}\frac{x^s}{y^j}+\frac{1}{2}\left[\frac{\mathrm{\Gamma }_\mu ^\nu }{x^k}\left(\frac{x^k}{y^j}\frac{x^r}{y^i}\frac{x^k}{y^i}\frac{x^r}{y^j}\right)e_r^\mu +\frac{f_i^\nu }{x^k}\frac{x^k}{y^j}\frac{f_j^\nu }{x^k}\frac{x^k}{y^i}\right]\hfill \end{array}$$ In a similar way (compare with ), it is easy to prove that the only vector fields of the form $$X=ϵ^i(x^j)\frac{}{x^i}+\left(\frac{ϵ^k}{x^q}e_k^\mu +D_\nu ^\mu (x^j)e_q^\nu +G_q^\mu (x^j)\right)\frac{}{e_q^\mu }$$ (12) where $`ϵ^i`$, $`D_\nu ^\mu `$ and $`G_q^\mu `$ are arbitrary local functions on $`M`$, can be $`𝒥`$prolonged to vector fields over $`𝒥()`$ as follows: $$𝒥(X)(z):=\rho _{\rho ^1(z)}(j_1(X))z𝒥()$$ (13) The resulting vector field has the form: $$𝒥(X)=ϵ^i(x^j)\frac{}{x^i}+\left(\frac{ϵ^k}{x^q}e_k^\mu +D_\nu ^\mu (x^j)e_q^\nu +G_q^\mu (x^j)\right)\frac{}{e_q^\mu }+\underset{i<j}{}h_{ij}^\mu \frac{}{E_{ij}^\mu }$$ where $$h_{ij}^\mu =\frac{1}{2}\left(\frac{D_\nu ^\mu }{x^j}e_i^\nu \frac{D_\nu ^\mu }{x^i}e_j^\nu +\frac{G_i^\mu }{x^j}\frac{G_j^\mu }{x^i}\right)+D_\nu ^\mu E_{ij}^\nu +\left(E_{ki}^\mu \frac{ϵ^k}{x^j}E_{kj}^\mu \frac{ϵ^k}{x^i}\right)$$ In the following discussion the central role will be played by a specific coordinate transformation in the space $`𝒥()`$. More precisely, the main idea consists in choosing the components of the spin–connections generated by the tetrads themselves as fiber coordinates on the bundle $`𝒥()`$. To see this point, let $`z=(x^i,e_i^\mu ,E_{ij}^\mu )`$ be an element of $`𝒥()`$, $`x=\widehat{\pi }(z)`$ its projection over $`M`$ and $`e^\mu `$ a representative tetrad belonging to the equivalence class $`z`$. Moreover, if $`g=\eta _{\mu \nu }e^\mu e^\nu `$ is the metric on $`M`$ induced by the tetrad $`e^\mu `$, denote by $`\mathrm{\Gamma }_{ih}^k`$ its associated Levi–Civita connection. The latter is a principal connection on $`L(M)`$ and can be pulled–back to a spin–connection $`\omega _{i\nu }^\mu `$ over $`P`$ by means of the tetrad $`e^\mu `$ itself (i.e. through the principal morphism $`i:PL(M)`$ associated to the tetrad $`e^\mu `$). The relation between the coefficients $`\mathrm{\Gamma }_{ih}^k`$ of the Levi–Civita connection and the coefficients $`\omega _{i\nu }^\mu `$ of the associated spin–connection, evaluated at the point $`x=\widehat{\pi }(z)M`$, is expressed by the equation $$\omega _{i\nu }^\mu (x)=e_k^\mu (x)\left(\mathrm{\Gamma }_{ij}^ke_\nu ^j(x)+\frac{e_\nu ^k(x)}{x^i}\right)$$ (14) In other and simpler words, the latter can be though as the Levi–Civita connection expressed in terms of the non–holonomic basis $`e^\mu (x)`$. If the coefficients $`\mathrm{\Gamma }_{ih}^k`$ are written in terms of the tetrad $`e^\mu `$ and its derivatives, one gets the well–known expression $$\omega _{i\nu }^\mu (x):=e_p^\mu (x)\left(\mathrm{\Sigma }_{ji}^p(x)\mathrm{\Sigma }_{ji}^p(x)+\mathrm{\Sigma }_{ij}^p(x)\right)e_\nu ^j(x)$$ (15) where $$\mathrm{\Sigma }_{ji}^p(x):=e_\lambda ^p(x)E_{ij}^\lambda (x)=e_\lambda ^p(x)\frac{1}{2}\left(\frac{e_i^\lambda (x)}{x^j}\frac{e_j^\lambda (x)}{x^i}\right)$$ (16) the Latin indexes being lowered and raised by means of the metric $`g=\eta _{\mu \nu }e^\mu e^\nu `$. Equations (15) and (16) show that the values of the coefficients of the spin–connection $`\omega _{i\nu }^\mu `$, evaluated in $`x=\widehat{\pi }(z)`$, are independent of the choice of the representative $`e^\mu `$ in the equivalence class $`z𝒥()`$. Moreover, the torsion–free condition for the connection $`\omega _{i\nu }^\mu `$ gives a sort of inverse relation of eq. (15) in the form $$2E_{ij}^\mu (x)=\frac{e_i^\mu (x)}{x^j}\frac{e_j^\mu (x)}{x^i}=\omega _{i\nu }^\mu (x)e_j^\nu (x)\omega _{j\nu }^\mu (x)e_i^\nu (x)$$ (17) Because of the metric compatibility condition $`\omega _i^{\mu \nu }:=\omega _{i\sigma }^\mu \eta ^{\sigma \nu }=\omega _i^{\nu \mu }`$, there exists a one-to-one correspondence between the values of the antisymmetric part of the derivatives $`E_{ij}^\mu (x)=\frac{1}{2}\left(\frac{e_i^\mu (x)}{x^j}\frac{e_j^\mu (x)}{x^i}\right)`$ and the coefficients of the spin–connection $`\omega _i^{\mu \nu }(x)`$ in the point $`x=\widehat{\pi }(z)`$. The above considerations allow us to take the quantities $`\omega _i^{\mu \nu }`$ as fiber coordinates of the bundle $`𝒥()`$, looking at the relations (15) and (17) as coordinate changes in $`𝒥()`$. It is a well–known fact that the coordinate transformations (2) induce the following transformation laws for the spin–connection coefficients $`\omega _i^{\mu \nu }`$: $$\overline{\omega }_i^{\mu \nu }=\mathrm{\Lambda }_\sigma ^\mu (x)\mathrm{\Lambda }_\gamma ^\nu (x)\frac{x^j}{\overline{x}^i}\omega _j^{\sigma \gamma }\mathrm{\Lambda }_\sigma ^\eta (x)\frac{\mathrm{\Lambda }_\eta ^\mu (x)}{x^h}\frac{x^h}{\overline{x}^i}\eta ^{\sigma \nu }$$ (18) where $`\mathrm{\Lambda }_\sigma ^\nu :=\mathrm{\Lambda }_\beta ^\alpha \eta _{\alpha \sigma }\eta ^{\beta \nu }=\left(\mathrm{\Lambda }^1\right)_\sigma ^\nu `$. Let us now define the variational principle from which we shall deduce the field equations for General Relativity directly on the only manifold $`𝒥()`$. To this end, we first introduce the $`4`$-form on $`𝒥()`$ locally described as $$\mathrm{\Theta }:=\frac{1}{4}ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu e_p^\nu \left(d\omega _i^{\lambda \sigma }ds_j+\omega _{j\eta }^\lambda \omega _i^{\eta \sigma }ds\right)$$ (19) where $`ds:=dx^1\mathrm{}dx^4`$, $`ds_i:=\frac{}{x^i}\text{ }\text{ }ds`$ and $`ϵ`$ denotes the Levi–Civita permutation symbol. The following result holds true ###### Proposition 1 The form $`\mathrm{\Theta }`$ (19) is invariant under the coordinate transformations (2), (18) on the manifold $`𝒥()`$. Proof. It is a direct check, taking eqs. (2), (18) and the identities $$\begin{array}{c}\hfill \overline{\omega }_{j\eta }^\tau \overline{\omega }_{i\sigma }^\eta =\mathrm{\Lambda }_\alpha ^\tau \mathrm{\Lambda }_\sigma ^\beta \frac{x^k}{\overline{x}^j}\frac{x^h}{\overline{x}^i}\omega _{k\lambda }^\alpha \omega _{h\beta }^\lambda +\mathrm{\Lambda }_\alpha ^\tau \frac{\mathrm{\Lambda }_\sigma ^\beta }{x^h}\frac{x^h}{\overline{x}^i}\frac{x^k}{\overline{x}^j}\omega _{k\beta }^\alpha +\\ \hfill \frac{\mathrm{\Lambda }_\alpha ^\tau }{x^h}\frac{x^h}{\overline{x}^j}\mathrm{\Lambda }_\sigma ^\beta \frac{x^k}{\overline{x}^i}\omega _{k\beta }^\alpha \frac{\mathrm{\Lambda }_\eta ^\tau }{x^h}\frac{x^h}{\overline{x}^j}\frac{\mathrm{\Lambda }_\sigma ^\eta }{x^k}\frac{x^k}{\overline{x}^i}\end{array}$$ explicitly into account. Being the form $`\mathrm{\Theta }`$ a covariant geometrical object, it can be used to define a variational problem on the bundle $`𝒥()`$, consisting in the study of the stationarity conditions for the functional $$A(\gamma ):=_D\gamma ^{}(\mathrm{\Theta })$$ (20) for every section $`\gamma :DM𝒥()`$, $`D`$ compact domain. The procedure is well known: we take a vertical vector field $`X`$ (with respect to the fibration $`𝒥()M`$) into account and denote by $`\mathrm{\Phi }_\xi `$ its flow; then, we deform any given section $`\gamma :M𝒥()`$ along $`X`$ by setting $`\gamma _\xi :=\mathrm{\Phi }_\xi \gamma `$. We name first variation of $`A`$ at $`\gamma `$ in the direction $`X`$ the expression (see, for example, ) $$\frac{\delta A}{\delta X}(\gamma ):=\frac{d}{d\xi }_D\gamma _\xi ^{}(\mathrm{\Theta })_{\|_{\xi =0}}=_D\gamma ^{}(X\text{ }\text{ }d\mathrm{\Theta })+_D\gamma ^{}(X\text{ }\text{ }\mathrm{\Theta })$$ (21) Finally, we look for sections $`\gamma :x(x^i,e_i^\mu (x),\omega _i^{\mu \nu }(x))`$ (critical points) obeying the ansatz $`\frac{\delta A}{\delta X}(\gamma )=0`$, for all compact domains $`D`$ and all infinitesimal deformations $`X`$ vanishing on the boundary $`D`$. According to eq. (21) and to the imposed boundary condition, a section $`\gamma `$ is critical if and only if it satisfies the equation $$\gamma ^{}(X\text{ }\text{ }\mathrm{\Theta })=0$$ (22) for every vector field $`X=X_i^\mu \frac{}{e_i^\mu }+\frac{1}{2}X_i^{\mu \nu }\frac{}{\omega _i^{\mu \nu }}`$ on $`𝒥()`$ (with $`X_i^{\mu \nu }=X_i^{\nu \mu }`$ when $`\mu >\nu `$). In order to make eq. (22) explicit, we calculate the differential of the form $`\mathrm{\Theta }`$, namely $$d\mathrm{\Theta }=\frac{1}{2}ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu de_p^\nu \left(d\omega _i^{\lambda \sigma }ds_j+\omega _{j\eta }^\lambda \omega _i^{\eta \sigma }ds\right)+\frac{1}{2}ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu e_p^\nu \omega _{j\eta }^\lambda d\omega _i^{\eta \sigma }ds$$ (23) ###### Proposition 2 The following identities $$ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu e_p^\nu \omega _{j\eta }^\lambda d\omega _i^{\eta \sigma }=ϵ^{qpij}ϵ_{\mu \rho \lambda \sigma }e_q^\mu e_p^\nu \omega _{j\nu }^\rho d\omega _i^{\lambda \sigma }$$ (24) hold true. Proof. Observing that the expressions $`ϵ^{qpij}e_q^\mu e_p^\nu `$ are antisymmetric in the indexes $`\mu `$ and $`\nu `$, the identities (24) will be proved if we can show that the antisymmetric combinations (still in the indexes $`\mu `$ and $`\nu `$) of the $`1`$-forms $`ϵ_{\mu \nu \lambda \sigma }\omega _{j\eta }^\lambda d\omega _i^{\eta \sigma }`$ and $`ϵ_{\mu \rho \lambda \sigma }\omega _{j\nu }^\rho d\omega _i^{\lambda \sigma }`$ coincide. In turn, the last assertion is mathematically equivalent to the fact that the following identities $$ϵ^{\mu \nu \alpha \beta }ϵ_{\mu \nu \lambda \sigma }\omega _{j\eta }^\lambda d\omega _i^{\eta \sigma }=ϵ^{\mu \nu \alpha \beta }ϵ_{\mu \rho \lambda \sigma }\omega _{j\nu }^\rho d\omega _i^{\lambda \sigma }$$ (25) hold true. Due to the traceless property $`\omega _{i\mu }^\mu =0`$, a direct calculation shows that both left and right hand sides of (25) are actually equal to $`2\left(\omega _{j\eta }^\alpha d\omega _i^{\eta \beta }\omega _{j\eta }^\beta d\omega _i^{\eta \alpha }\right)`$. Making use of the identities (24), we can rewrite the expression (23) in the form $$d\mathrm{\Theta }=\frac{1}{2}ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu de_p^\nu \left(d\omega _i^{\lambda \sigma }ds_j+\omega _{j\eta }^\lambda \omega _i^{\eta \sigma }ds\right)\frac{1}{2}ϵ^{qpij}ϵ_{\mu \rho \lambda \sigma }e_q^\mu e_p^\nu \omega _{j\nu }^\rho d\omega _i^{\lambda \sigma }ds$$ (26) Now, given a vector field $`X=X_i^\mu \frac{}{e_i^\mu }+\frac{1}{2}X_i^{\mu \nu }\frac{}{\omega _i^{\mu \nu }}`$ on $`𝒥()`$, we easily have $$\begin{array}{c}\hfill X\text{ }\text{ }d\mathrm{\Theta }=\frac{1}{2}ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu \left(d\omega _i^{\lambda \sigma }ds_j+\omega _{j\eta }^\lambda \omega _i^{\eta \sigma }ds\right)X_p^\nu +\\ \hfill \frac{1}{2}ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu \left(de_p^\nu ds_j+e_p^\rho \omega _{j\rho }^\nu ds\right)X_i^{\lambda \sigma }\end{array}$$ (27) In conclusion, the imposition of condition (22) yields two sets of final equations $$ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu \left(\frac{e_p^\nu }{x^j}+\omega _{j\rho }^\nu e_p^\rho \right)=0$$ (28a) $$\frac{1}{2}ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu \left(\frac{\omega _i^{\lambda \sigma }}{x^j}+\omega _{j\eta }^\lambda \omega _i^{\eta \sigma }\right)=0$$ (28b) clearly equivalent to Einstein equations (provided that $`\mathrm{det}(e_i^\mu )0`$). Indeed, eqs. (28a) ensure the kinematic admissibility (the holonomy) of the critical section $`\gamma `$, namely $$2E_{pj}^\nu (x)=\omega _{p\rho }^\nu (x)e_j^\rho (x)\omega _{j\rho }^\nu (x)e_p^\rho (x)=\frac{e_p^\nu }{x^j}(x)\frac{e_j^\nu }{x^p}(x)$$ so that the quantities $`\omega _{i\nu }^\mu (x)`$ identify with the coefficients of the spin connection associated to the Levi–Civita connection induced by the metric $`g_{ij}(x)=\eta _{\mu \nu }e_i^\mu (x)e_j^\nu (x)`$. Therefore eqs. (28b) are identical to $$\frac{1}{4}ϵ^{qpij}ϵ_{\mu \nu \lambda \sigma }e_q^\mu (x)R_{ji}^{\lambda \sigma }(x)=0$$ $`R_{ji}^{\lambda \sigma }(x):=\frac{\omega _i^{\lambda \sigma }}{x^j}(x)\frac{\omega _j^{\lambda \sigma }}{x^i}(x)+\omega _{j\eta }^\lambda (x)\omega _i^{\eta \sigma }(x)\omega _{i\eta }^\lambda (x)\omega _j^{\eta \sigma }(x)`$ denoting the curvature tensor of the metric $`g`$. It is worth noticing that the restriction regarding the verticality of the infinitesimal deformations $`X`$ can be removed, since condition (22) automatically implies $`\gamma ^{}(X\text{ }\text{ }\mathrm{\Theta })=0`$ $`XD^1(𝒥())`$. This last fact is important in order to extend the study of Noether vector fields, conserved currents and symmetries to the present geometrical setting. In particular, any given vector field $`Z`$ on $`𝒥()`$ will be called a Noether vector field if it satisfies the ansatz $$L_Z\mathrm{\Theta }=\omega +d\alpha $$ (29) where $`\omega `$ is a $`4`$-form belonging to the ideal generated by the contact forms and $`\alpha `$ is any $`3`$-form on $`𝒥()`$. If $`Z`$ satisfies the trivial case $`L_Z\mathrm{\Theta }_L=0`$ and projects to $`M`$, then $`Z`$ is an infinitesimal dynamical symmetry (namely, its flow drags critical sections into as many critical sections). It is also easy to verify that whenever a Noether vector field $`Z`$ is a $`𝒥`$-prolongation of some vector field (12) on $``$, it again results into an infinitesimal dynamical symmetry. Moreover, a corresponding conserved current is always associated with any Noether vector field $`Z`$ . In fact, given a critical section $`\gamma `$, one has $$d\gamma ^{}\left(Z\text{ }\text{ }\mathrm{\Theta }\alpha \right)=\gamma ^{}\left(\omega Z\text{ }\text{ }d\mathrm{\Theta }\right)=0$$ (30) The current $`\gamma ^{}\left(Z\text{ }\text{ }\mathrm{\Theta }\alpha \right)`$ is then conserved on shell. We conclude this letter by noticing that a new geometrical description of the combined theory of gravitation and Yang–Mills fields within the framework of $`𝒥`$-bundles can be obtained, joining the present geometrical approach with the one developed in . The matter is straightforward and follows the lines already illustrated in for the tetrad–affine formulation; for brevity reasons, we leave the details to the reader.
warning/0506/cond-mat0506565.html	ar5iv	text	# Thermal broadening of the J-band in disordered linear molecular aggregates: A theoretical study ## I Introduction Ever since the discovery of aggregation of cyanine dye molecules by Jelley Jelley36 and Scheibe Scheibe36 the width of the absorption band of these linear aggregates (the J-band) has attracted much attention. At low temperatures, the J-band may be as narrow as a few tens of cm<sup>-1</sup>’s (for pseudoisocyanine), while at room temperature it typically is a few hundred cm<sup>-1</sup>.Kobayashi96 The small width at low temperature is generally understood as resulting from the excitonic nature of the optical excitations,Franck38 which leads to exchange narrowing of the inhomogeneous broadening of the transitions of individual molecules.Knapp84 As the optically dominant exciton states in inhomogeneous J-aggregates occur below the exciton band edge, their vibration-induced dephasing is strongly suppressed at low temperatures. Thus, the J-band is inhomogeneously broadened, except for a small residual homogeneous component ($``$ 0.1 cm<sup>-1</sup>) caused by spontaneous emission of the individual exciton states underlying the spectrum. Upon increasing the temperature, the vibration-induced dephasing of the exciton states increases, the J-band broadens, and obtains a more homogeneous character. Over the past twenty years, the temperature dependence of the J-band width and the homogeneous broadening of the exciton states, has been studied by several authors. In 1987, De Boer, Vink, and Wiersma deBoer87 performed accumulated-echo experiments to study the temperature dependence of the pure dephasing time in J-aggregates of pseudo-isocyanine bromide (PIC-Br) and found that within the temperature range 1.5 K to 100 K, this time is linearly proportional to the occupation number of a mode with an energy of 9 cm<sup>-1</sup>. They assigned this to a librational mode of the aggregate. Later on, Fidder, Knoester, and Wiersma Fidder90 (see also Ref. Fidder91, and Fidder93, ) showed that a similar study carried our over a wider temperature range, 1.5 K to 190 K, required the occupation numbers of three vibration modes, at 9 cm<sup>-1</sup>, 305 cm<sup>-1</sup>, and 973 cm<sup>-1</sup>. Using the hole burning technique, Hirschmann and Friedrich Hirschmann89 studied the homogeneous width of the exciton states in pseudoisocyanine iodide (PIC-I) over the temperature range 350 mK to 80 K. They were able to fit their measurements by a superposition of two exponentials, with activation energies 27 cm<sup>-1</sup> and 330 cm<sup>-1</sup>, and they attributed the broadening to scattering of the excitons on an acoustic mode and an optical mode, respectively, of the aggregate. Finally, in 1997 Renge and Wild Renge97 measured the temperature dependent width, $`\mathrm{\Delta }(T)`$, of the total J-band of pseudoisocyanine chloride (PIC-Cl) and fluoride (PIC-F) over the wide temperature range 10 K to 300 K. They found that over this entire range their data closely obeyed a power-law dependence $`\mathrm{\Delta }(T)=\mathrm{\Delta }(0)+bT^p`$ with the exponent $`p=3.4`$. They suggested the scattering of the excitons on host vibrations as a possible source of this behavior. The above overview clearly demonstrates that the temperature dependence of the most important characteristic of J-aggregates, namely the J-band, is not understood. It even is not clear what is the source of thermal broadening: dephasing due to vibrational modes of the aggregate itself or due to modes of the host matrix. Theoretically, the study of the exciton dephasing in J-aggregates is complicated by the fact that static disorder plays an important role in these systems, as is clear from the strongly asymmetric low-temperature lineshape. Fidder91 ; Fidder93 The simultaneous treatment of disorder, leading to exciton localization, and scattering on vibrational modes, is a problem that requires extensive numerical simulations. Such simulations have been used previously to model the optical response of polysilanes, Shimitsu98 ; Shimizu01 the temperature dependent fluorescence Bednarz03 and transport Malyshev03 properties of J-aggregates, and single molecular spectroscopy of circular aggregates. Dempster01 Alternatively, stochastic models have been applied to describe the scattering of excitons on vibrations in disordered aggregates and the resulting spectral line shape Barvik99 ; Warns03 and relaxation dynamics. Lemaistre99 In this paper, we report on a systematic theoretical study of dephasing of weakly localized Frenkel excitons in one-dimensional systems, focusing on the effect of the scattering of the excitons on the vibrational modes (phonons) of the host. For a chain-like configuration it seems physically reasonable to assume that the coupling to host vibrations dominates the dynamics of the excitons. For coupling to vibrations in the chain, one expects self-trapped exciton states, Agranovich99 for which in most aggregates no clear signature is found. We describe the phonons by a Debye model and consider one-phonon as well as (elastic and inelastic) two-phonon contributions to the dephasing. The scattering rates between the various, numerically obtained, exciton states are derived using the Fermi Golden Rule. Due to the disorder, the dephasing rates of individual exciton states are distributed over a wide range, in particular at low temperature, making it meaningless to associate the width of individual exciton levels with the J-band width. Rather, this width is determined by direct simulation of the total absorption spectrum and is found to scale with temperature according to a power law. The outline of this paper is as follows: In Sec. II we introduce the Hamiltonian for excitons in a disordered chain and coupled to host vibrations. General expressions for the exciton dephasing rates are derived in Sec. III. Section IV deals with the temperature dependence of the dephasing rates for the disorder-free case, where analytical expressions can be obtained. In Sec. V we give the results of our numerical simulations for the dephasing rates in the presence of disorder and analyze their temperature dependence and fluctuations as well as the total J-band width. We compare to experiment in Sec. VI and discuss an alternative mechanism of dephasing due to the coupling of excitons to a local vibration. In Section VII we present our conclusions. ## II Model We consider an ensemble of J-aggregates embedded in a disordered host matrix. The aggregates are assumed to be decoupled from each other, while they interact with the host. A single aggregate is modeled as an open linear chain of $`N`$ coupled two-level monomers with parallel transition dipoles. The interaction between a particular monomer and the surrounding host molecules in the equilibrium configuration leads to shifts in the monomer’s transition energy. Due to the host’s structural disorder, this shift is different for each monomer in the aggregate, giving rise to on-site (diagonal) disorder. Moreover, vibrations of the host couple to the aggregate excited states, because the associated displacements away from the equilibrium configuration dynamically affect the monomer transition energies. Accounting for these shifts up to second order in the molecular displacements, the resulting Hamiltonian in the site representation reads: $$H=H^{\mathrm{ex}}+H^{\mathrm{bath}}+V^{(1)}+V^{(2)},$$ (1a) with $$H^{\mathrm{ex}}=\underset{n=1}{\overset{N}{}}\epsilon _n\|nn\|+\underset{n,m=1}{\overset{N}{}}J_{nm}\|nm\|,$$ (1b) $$H^{\mathrm{bath}}=\underset{q}{}\omega _qa_q^{}a_q,$$ (1c) $$V^{(1)}=\underset{n=1}{\overset{N}{}}\underset{q}{}V_{nq}^{(1)}\|nn\|(a_q+a_q^{}),$$ (1d) $$V^{(2)}=\underset{n=1}{\overset{N}{}}\underset{qq^{}}{}V_{nqq^{}}^{(2)}\|nn\|(a_q+a_q^{})(a_q^{}+a_q^{}^{}).$$ (1e) Here, $`H^{\mathrm{ex}}`$ is the bare Frenkel exciton Hamiltonian, with $`\|n`$ denoting the state in which the $`n`$th monomer is excited and all the other monomers are in the ground state. The monomer excitation energies, $`\epsilon _1,\epsilon _2,\mathrm{},\epsilon _N`$, are uncorrelated stochastic gaussian variables, with mean $`\overline{\epsilon }`$ and standard deviation $`\sigma `$, referred to as the disorder strength. Hereafter, $`\overline{\epsilon }`$ is set to zero. The resonant interactions $`J_{nm}`$ are considered to be nonrandom and are assumed to be of dipolar origin: $`J_{nm}=J/\|nm\|^3`$$`(J_{nn}0)`$, with $`J>0`$ denoting the nearest-neighbor coupling. $`H^{\mathrm{bath}}`$ describes the vibrational modes of the host, labeled $`q`$ and with the spectrum $`\omega _q`$ ($`\mathrm{}=1`$). The operator $`a_q(a_q^{})`$ annihilates (creates) a vibrational quantum in mode $`q`$. Finally, the operators $`V^{(1)}`$ and $`V^{(2)}`$ describe the linear and quadratic exciton-vibration coupling, respectively, where the quantities $`V_{nq}^{(1)}`$ and $`V_{nqq^{}}^{(2)}`$ indicate their strengths. We do not derive explicit expressions for these coupling strengths, as we aim to treat them on a phenomenological basis. In particular, owing to the disordered nature of the host, we consider these strengths stochastic quantities, for which we only specify the following stochastic properties with respect to the site index $`n`$: $$V_{nq}^{(1)}=V_{nqq^{}}^{(2)}=0,$$ (2a) $$V_{mq}^{(1)}V_{nq}^{(1)}=\delta _{mn}\left\|V_q^{(1)}\right\|^2,$$ (2b) $$V_{mqq^{}}^{(2)}V_{nqq^{}}^{(2)}=\delta _{mn}\left\|V_{qq^{}}^{(2)}\right\|^2,$$ (2c) where the angular brackets denote averaging over realizations of $`V_{nq}^{(1)}`$ and $`V_{nqq^{}}^{(2)}`$. The Kronecker symbol in Eqs. (2b) and (2c) implies that the surroundings of different monomers in the aggregate are not correlated. In general, the Hamiltonian Eq. (1) can not be diagonalized analytically. If the exciton-vibration coupling is not too strong, the method of choice is first to find the exciton eigenstates by numerical diagonalization and then to consider the scattering on the basis of these eigenstates. Explicitly, the exciton states follow from the eigenvalue problem $$\underset{m=1}{\overset{N}{}}H_{nm}^{\mathrm{ex}}\phi _{\nu m}=E_\nu \phi _{\nu n},\nu =1,2,\mathrm{},N,$$ (3a) where $`H_{nm}^{\mathrm{ex}}=n\|H^{\mathrm{ex}}\|m`$ and $`E_\nu `$ is the eigenenergy of the exciton state $`\|\nu `$: $$\|\nu =\underset{n=1}{\overset{N}{}}\phi _{\nu n}\|n.$$ (3b) In the exciton representation, the exciton-vibration interactions take the form $$V^{(1)}=\underset{\mu ,\nu =1}{\overset{N}{}}\underset{q}{}V_{\mu \nu q}^{(1)}\|\mu \nu \|(a_q+a_q^{}),$$ (4a) $$V^{(2)}=\underset{\mu ,\nu =1}{\overset{N}{}}\underset{qq^{}}{}V_{\mu \nu qq^{}}^{(2)}\|\mu \nu \|(a_q+a_q^{})(a_q^{}+a_q^{}^{}),$$ (4b) where $`V_{\mu \nu q}^{(1)}`$ and $`V_{\mu \nu qq^{}}^{(2)}`$ are the matrix elements of the vibration-induced scattering of an exciton from state $`\|\nu `$ to state $`\|\mu `$, given by $$V_{\mu \nu q}^{(1)}=\underset{n=1}{\overset{N}{}}V_{nq}^{(1)}\phi _{\mu n}\phi _{\nu n},$$ (5a) $$V_{\mu \nu qq^{}}^{(2)}=\underset{n=1}{\overset{N}{}}V_{nqq^{}}^{(2)}\phi _{\mu n}\phi _{\nu n}.$$ (5b) Spectroscopic data on J-aggregates clearly reveal that the exciton-vibration coupling in these systems is usually weak. For the prototypical J-aggregates of PIC, this claim is corroborated by two facts: (i) - the narrowness of the J-band, which only is a few tens of cm<sup>-1</sup> at liquid helium temperature and becomes several times broader at room temperature, and (ii) - the absence of a fluorescence Stokes shift of the J-band (see, e.g., Ref. Fidder90, ). The extended nature of the exciton states in J-aggregates helps to reduce the exciton-vibration coupling, as it leads to averaging of the static as well as dynamic fluctuations of the site energies, effects known as exchange Knapp84 and motional Wubs98 ; Malyshev98 narrowing, respectively. The weakness of the exciton-vibration coupling allows one to calculate the scattering and dephasing rates of the excitons through perturbation theory. This analysis is presented in the next section. ## III Dephasing rates Following the arguments given at the end of Sec. II, we will use Fermi’s Golden Rule to calculate the rate for scattering of excitons from one localized state, $`\|\nu `$, to another one, $`\|\mu `$. The result reads $`W_{\mu \nu }^{(\xi )}`$ $`=`$ $`2\pi {\displaystyle \underset{f}{}}{\displaystyle \underset{i}{}}\rho (\mathrm{\Omega }_i)\left\|\mu ,f\|V^{(\xi )}\|\nu ,i\right\|^2`$ (6) $`\times `$ $`\delta (E_\mu E_\nu +\mathrm{\Omega }_f\mathrm{\Omega }_i).`$ Here, the superscript $`\xi =1,2`$ distinguishes between one- and two-vibration-assisted exciton scattering. Furthermore, $`\mathrm{\Omega }_i`$ and $`\mathrm{\Omega }_f`$ are the energies of the vibration bath in the initial ($`\|i=\|\{n_q\}_i`$) and final ($`\|f=\|\{n_q\}_f`$) states, respectively, where $`\{n_q\}`$ denotes the set of occupation numbers of the vibrational modes. The quantity $`\rho (\mathrm{\Omega }_i)`$ is the equilibrium density matrix of the initial state of the bath. Finally, the angular brackets indicate that we average over the stochastic realizations of the surroundings of each monomer in the aggregate. ### III.1 Linear exciton-vibration coupling In a one-phonon-assisted scattering process, the occupation number of one phonon mode $`q`$ increases or decreases by one, corresponding to emission and absorption of a vibrational quantum, respectively. Consequently, $`\mathrm{\Omega }_f\mathrm{\Omega }_i=\pm \omega _q`$. Substituting the explicit form of the operator $`V^{(1)}`$ from Eq. (1d) into Eq. (6), we obtain $`W_{\mu \nu }^{(1)}`$ $`=`$ $`2\pi {\displaystyle \underset{n=1}{\overset{N}{}}}\phi _{\mu n}^2\phi _{\nu n}^2{\displaystyle \underset{q}{}}\left\|V_q^{(1)}\right\|^2`$ (7) $`\times `$ $`[[\overline{n}(\omega _q)+1]\delta (\omega _{\mu \nu }+\omega _q)`$ $`+`$ $`\overline{n}(\omega _q)\delta (\omega _{\mu \nu }\omega _q)],`$ where $`\overline{n}(\omega _q)=[\mathrm{exp}(\omega _q/T)1]^1`$ is the mean occupation number of the vibrational mode $`q`$ (the Boltzmann constant $`k_B=1`$) and $`\omega _{\mu \nu }=E_\mu E_\nu `$. In deriving Eq. (7), we used the properties of the stochastic function $`V_{nq}^{(1)}`$ given by Eqs. (2a) and (2b). Defining the one-vibration spectral density as $$^{(1)}(\omega )2\pi \underset{q}{}\left\|V_q^{(1)}\right\|^2\delta (\omega \omega _q),$$ (8) we can rewrite Eq. (7) in the form $`W_{\mu \nu }^{(1)}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{N}{}}}\phi _{\mu n}^2\phi _{\nu n}^2^{(1)}(\|\omega _{\mu \nu }\|)`$ (12) $`\times `$ $`\{\begin{array}{cc}\overline{n}(\omega _{\mu \nu }),\hfill & \hfill \omega _{\mu \nu }>0,\\ & \\ \overline{n}(\omega _{\mu \nu })+1,\hfill & \hfill \omega _{\mu \nu }<0.\end{array}`$ As we observe, the rate $`W_{\mu \nu }^{(1)}`$ is proportional to the overlap integral of the site occupation probabilities, $`\phi _{\mu n}^2`$ and $`\phi _{\nu n}^2`$, of the exciton states involved. First of all, this leads to a strong suppression of the scattering rate if states $`\|\mu `$ and $`\|\nu `$ overlap weakly or not at all. Second, as the low-energy exciton states in a disordered chain exhibit large fluctuations in their localization size,Malyshev01 also the scattering rates may undergo large fluctuations (see Sec. V). The dependence of $`W_{\mu \nu }^{(1)}`$ on the energy mismatch $`\omega _{\mu \nu }`$ is determined by the one-phonon spectral density $`^{(1)}(\omega )`$. Characterizing this function requires knowledge of the vibrational spectrum $`\omega _q`$ as well as the $`q`$ dependence of the exciton-vibration coupling $`V_q^{(1)}`$. For the special case of scattering on acoustic phonons of a relatively long wavelength, we have $`^{(1)}(\omega )\omega ^3`$. This behavior results from the $`\omega ^2`$-dependence of the density of states of acoustic phonons, combined with the fact that in the long-wavelength limit $`\|V_q^{(1)}\|^2\omega _q`$.Davydov71 ; Bednarz02 One may consider this a Debye-like model, in which one replaces the summation over the mode index $`q`$ by an integration over the frequency $`\omega _q`$ according to the well-known rule: $$\underset{q}{}C_0^{\omega _c}𝑑\omega _q\omega _q^2.$$ (13) Here, $`C`$ is an irrelevant constant, which we will incorporate in an overall free parameter (see below), and $`\omega _c`$ is a cutoff frequency. It is important to note that $`\omega _c`$ is not necessarily related to the Debye frequency: the generally very complex density of vibrational states in a disordered solid may on average exhibit an $`\omega ^2`$ scaling up to a given frequency $`\omega _c`$. Inspired by the above, we consider a slightly wider class of one-phonon spectral densities, given by $$^{(1)}(\omega )=W_0^{(1)}\left(\frac{\omega }{J}\right)^\alpha \mathrm{\Theta }(\omega _c\omega ).$$ (14) Here, $`W_0^{(1)}`$ is a free parameter in the model, which characterizes the overall strength of the one-vibration-assisted scattering rates. It absorbs a number of other parameters characteristic for the host lattice (such as the velocity of sound), the constant $`C`$ from Eq. (13), as well as the strength of the transfer interaction $`J`$ (for details, see Ref. Bednarz02, ). $`\mathrm{\Theta }(x)`$ is the Heaviside step function. When performing numerical simulations, we will mostly use $`\alpha =3`$, for which the spectral density of acoustic phonons in the long-wave limit is recovered. In some instances, however, we will discuss how results depend on the exponent $`\alpha `$. Several observations support considering a Debye-like vibration spectral density, even for a disordered host. Thus, for strongly disordered Yb<sup>3+</sup> doped phosphate glasses, a parabolic behavior of the one-phonon spectral density was found over a rather wide range of measurement (0 to 100 cm<sup>-1</sup>). Basiev87 Furthermore, closely related spectral densities of the form $`^{(1)}(\omega )(\omega /\omega _c)^\alpha \mathrm{exp}(\omega /\omega _c)`$ (or linear combinations of such functions) have been used successfully to fit the optical dynamics in photosynthetic antenna complexes (see, e.g., Refs. Kuhn97, ; May00, ; Renger01, ; Brueggemann04, ). One-phonon-assisted scattering results in the transition of an exciton from a given state $`\|\nu `$ to state $`\|\mu `$, where necessarily $`\mu \nu `$. In other words, this type of scattering changes the occupation probabilities of the exciton states and thus causes population (or energy) relaxation. The population relaxation in turn contributes to the dephasing of state $`\|\nu `$. The corresponding dephasing rate is given by (see, e.g., Ref. Blum96, ): $$\mathrm{\Gamma }_\nu ^{(1)}\frac{1}{2}\underset{\mu (\nu )}{}W_{\mu \nu }^{(1)}.$$ (15) Thus, $`\mathrm{\Gamma }_\nu ^{(1)}`$ represents the one-phonon-assisted contribution to the homogeneous broadening of the excitonic level $`\nu `$. We note that the $`\mathrm{\Gamma }_\nu ^{(1)}`$ indirectly depend on temperature through the $`\overline{n}(\omega _{\mu \nu })`$ \[cf. Eq. (12)\]. The temperature dependence of the sum over scattering rates in Eq. (15) and the corresponding width of the total exciton absorption spectrum will be analyzed in Secs. IV and V. ### III.2 Quadratic exciton-phonon coupling When excitons scatter on the second-order displacements of the host molecules, described by the operator $`V^{(2)}`$, the occupation numbers of two phonon modes $`q`$ and $`q^{}`$ with frequencies $`\omega _q`$ and $`\omega _q^{}`$ change by $`\pm 1`$. Thus, $`\mathrm{\Omega }_f\mathrm{\Omega }_i=\pm \omega _q\pm \omega _q^{}`$, where any combination of plus and minus is allowed. The corresponding scattering rates $`W_{\mu \nu }^{(2)}`$ are obtained from Eq. (6), taking into account the stochastic properties of $`V_{nqq^{}}^{(2)}`$ given by Eqs. (2a) and (2c): $`W_{\mu \nu }^{(2)}`$ $`=`$ $`2\pi {\displaystyle \underset{n=1}{\overset{N}{}}}\phi _{\mu n}^2\phi _{\nu n}^2{\displaystyle \underset{qq^{}}{}}\left\|V_{qq^{}}^{(2)}\right\|^2`$ (16) $`\times `$ $`[[\overline{n}(\omega _q)+1][\overline{n}(\omega _q^{})+1]\delta (\omega _{\mu \nu }+\omega _q+\omega _q^{})`$ $`+`$ $`2\overline{n}(\omega _q)\left[\overline{n}(\omega _q^{})+1\right]\delta (\omega _{\mu \nu }\omega _q+\omega _q^{})`$ $`+`$ $`\overline{n}(\omega _q)\overline{n}(\omega _q^{})\delta (\omega _{\mu \nu }\omega _q\omega _q^{})].`$ If in analogy to the one-vibration-assisted scattering, we define the two-vibration spectral density $`^{(2)}(\omega ,\omega ^{})`$ as $$^{(2)}(\omega ,\omega ^{})2\pi \underset{qq^{}}{}\left\|V_{qq^{}}^{(2)}\right\|^2\delta (\omega \omega _q)\delta (\omega ^{}\omega _q^{}),$$ (17) the scattering rate $`W_{\mu \nu }^{(2)}`$ takes the form $`W_{\mu \nu }^{(2)}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{N}{}}}\phi _{\mu n}^2\phi _{\nu n}^2{\displaystyle d\omega d\omega ^{}^{(2)}(\omega ,\omega ^{})}`$ (18) $`\times `$ $`[[\overline{n}(\omega )+1][\overline{n}(\omega ^{})+1]\delta (\omega _{\mu \nu }+\omega +\omega ^{})`$ $`+`$ $`2\overline{n}(\omega )\left[\overline{n}(\omega ^{})+1\right]\delta (\omega _{\mu \nu }\omega +\omega ^{})`$ $`+`$ $`\overline{n}(\omega )\overline{n}(\omega ^{})\delta (\omega _{\mu \nu }\omega \omega ^{})].`$ Similar to $`^{(1)}(\omega )`$ \[Eq. (8)\], we will use a parametrization $$^{(2)}(\omega ,\omega ^{})=\frac{W_0^{(2)}}{J}\left(\frac{\omega \omega ^{}}{J^2}\right)^\alpha \mathrm{\Theta }(\omega _c\omega )\mathrm{\Theta }(\omega _c\omega ^{}),$$ (19) where $`W_0^{(2)}`$ is a free parameter that characterizes the overall strength of the two-vibration-assisted scattering rates \[cf. $`W_0^{(1)}`$ and the discussion following Eq. (14)\]. Two main types of two-phonon-assisted processes may be distinguished. Similarly to the one-phonon case, an inelastic channel exists, where scattering occurs between different exciton states, thus giving rise to population relaxation. However, also an elastic channel is present, in which an exciton is scattered by emitting and absorbing a phonon of the same energy, and the final exciton state is identical to the initial one. This process results in pure dephasing of the exciton state, with a rate given by $`W_{\nu \nu }^{(2)}=2{\displaystyle \underset{n=1}{\overset{N}{}}}\phi _{\nu n}^4{\displaystyle d\omega ^{(2)}(\omega ,\omega )\overline{n}(\omega )\left[\overline{n}(\omega )+1\right]}.`$ (20) The quantity $`_{n=1}^N\phi _{\nu n}^4`$ is recognized as the inverse participation ratio, Thouless74 which is inversely proportional to the localization size of the exciton state $`\|\nu `$. Thus, we see that pure dephasing is suppressed for more extended states, an effect known as the motional narrowing. Wubs98 ; Malyshev98 Like in the one-phonon assisted process, the scattering rate between different states $`\|\nu `$ and $`\|\mu `$ is proportional to the overlap of their site occupations. The final expressions for the rates $`W_{\mu \nu }^{(2)}`$ resulting from Eqs. (18) and (19) are derived in Appendix A; they depend on the sign of $`\omega _{\mu \nu }`$ as well as on the relation between $`\|\omega _{\mu \nu }\|`$ and $`\omega _c`$. Distinction is made between three inelastic channels: downward ($``$), in which two phonons are emitted, cross ($``$), in which one phonon is absorbed and another is emitted, and upward ($``$), in which two phonons are absorbed. The fourth type of scattering is the elastic (pure dephasing) channel, discussed above already. When calculating the two-phonon assisted dephasing rate $`\mathrm{\Gamma }_\nu ^{(2)}`$ of state $`\|\nu `$, we will account for elastic as well as inelastic contributions: $$\mathrm{\Gamma }_\nu ^{(2)}=\frac{1}{2}\left[W_{\nu \nu }^{(2)}+\underset{\mu (\nu )}{}W_{\mu \nu }^{(2)}\right].$$ (21) This rate depends on temperature as a result of the mean occupation numbers $`\overline{n}(\omega )`$ and $`\overline{n}(\omega ^{})`$ of the vibrational modes. ## IV Disorder-free aggregate In order to gain insight in the temperature dependence of the dephasing rates and the absorption bandwidth, it is useful to start by considering a homogeneous aggregate, i.e., $`\sigma =0`$. Analytical results can then be obtained if we restrict the resonant interactions $`J_{nm}`$ to nearest-neighbor ones. In this approximation (which we will relax in our numerical analysis), we have $$\phi _{\nu n}=\left(\frac{2}{N+1}\right)^{1/2}\mathrm{sin}\frac{\pi \nu n}{N+1},$$ (22a) $$E_\nu =2J\mathrm{cos}\frac{\pi \nu }{N+1}.$$ (22b) The corresponding overlap integrals occurring in Eqs. (7) and (16) now read $$\underset{n=1}{\overset{N}{}}\phi _{\mu n}^2\phi _{\nu n}^2=\frac{1}{N+1}\left[1+\frac{1}{2}\left(\delta _{\mu \nu }+\delta _{\mu +\nu ,N+1}\right)\right].$$ (23) ### IV.1 One-phonon-assisted dephasing In a homogeneous linear chain, the lowest exciton state, $`\|\nu =1`$, contains almost all oscillator strength, thus dominating the absorption spectrum.Knapp84 ; Fidder90 Therefore, the dephasing rate of this state, $`\mathrm{\Gamma }_1^{(1)}`$, is of primary interest. As follows from Eq. (15), it is determined by the sum over scattering rates to the other exciton states ($`\mu 1`$), all of which are higher in energy. In order to evaluate $`\mathrm{\Gamma }_1^{(1)}`$, we replace the summation in Eq. (15) by an integration, $`\left(_\mu [(N+1)/\pi ]𝑑K\right)`$, which is allowed if $`N1`$ and $`TE_2E_1`$. We will also assume that $`TJ`$, which implies that the relevant exciton levels are those near the lower exciton band edge, where $`E_\mu =2J+JK^2`$ with $`K=\pi \mu /(N+1)`$. Changing the integration variable to $`x=JK^2/T`$, using $`E_\mu E_1JK^2`$ and replacing the lower integration limit $`\pi /(N+1)`$ by zero, we obtain $$\mathrm{\Gamma }_1^{(1)}=\frac{W_0^{(1)}}{4\pi }\left(\frac{T}{J}\right)^{\alpha +\frac{1}{2}}_0^{\omega _c/T}𝑑x\frac{x^{\alpha \frac{1}{2}}}{e^x1}.$$ (24) For temperatures $`T\omega _c`$, the upper integration limit may be extended to infinity, and we arrive at $$\mathrm{\Gamma }_1^{(1)}=\frac{W_0^{(1)}}{4\pi }\mathrm{\Gamma }\left(\alpha +\frac{1}{2}\right)\zeta \left(\alpha +\frac{1}{2}\right)\left(\frac{T}{J}\right)^{\alpha +\frac{1}{2}},$$ (25) where $`\mathrm{\Gamma }(z)`$ and $`\zeta (z)`$ are the gamma-function and the Rieman zeta-function, respectively. Thus, for $`T\omega _c`$ the one-phonon-assisted dephasing rate shows a power-law temperature dependence. Note that for our model of acoustic phonons ($`\alpha =3`$), $`\mathrm{\Gamma }_1^{(1)}`$ increases quite steeply, namely as $`T^{7/2}`$. From numerical evaluation of $`\mathrm{\Gamma }_1^{(1)}`$ for a homogeneous chain with all dipole-dipole interactions, we have found that the exponent 7/2 is increased to 3.85, mainly as a consequence of logarithmic corrections in the exciton dispersion near the lower band edge.Malyshev95 ; Didraga04 If we go beyond the parabolic range of the energy spectrum, the growth becomes even steeper; the exponent then tends to 4, because the density of states becomes a constant towards the center of the band. In the opposite limit $`T\omega _c`$, the exponential in the denominator of Eq. (24) can be expanded in a Taylor series. Up to second order, one obtains $$\mathrm{\Gamma }_1^{(1)}=\frac{W_0^{(1)}}{2\pi (2\alpha 1)}\left(\frac{\omega _c}{J}\right)^{\alpha \frac{1}{2}}\frac{T}{J},$$ (26) which simply reflects the linear high-temperature dependence of the mean occupation number $`\overline{n}(\omega )`$. Obviously, this scaling also holds in the presence of disorder. ### IV.2 Pure dephasing We now turn to the temperature dependence of the pure dephasing rate $`\mathrm{\Gamma }_\nu ^{(2)}=(1/2)W_{\nu \nu }^{(2)}`$. As for a homogeneous aggregate this rate does not depend on the state index $`\nu `$, we will simply denote it as $`\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}`$. Using the explicit form of $`^{(2)}(\omega ,\omega ^{})`$ given by Eq. (20), we arrive at $$\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}=\frac{3}{2}\frac{W_0^{(2)}}{N+1}\left(\frac{T}{J}\right)^{2\alpha +1}_0^{\omega _c/T}dx\frac{x^{2\alpha }e^x}{(e^x1)^2}.$$ (27) From Eq. (27) it follows that for $`T\omega _c`$ ($`\omega _c/T\mathrm{}`$), $$\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}=\frac{3}{2}\frac{W_0^{(2)}}{N+1}\mathrm{\Gamma }(2\alpha +1)\zeta (2\alpha )\left(\frac{T}{J}\right)^{2\alpha +1},$$ (28a) while for $`T\omega _c`$ ($`\omega _c/T0`$), $$\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}=\frac{3}{2}\frac{W_0^{(2)}}{(2\alpha 1)(N+1)}\left(\frac{\omega _c}{J}\right)^{2\alpha 1}\left(\frac{T}{J}\right)^2.$$ (28b) For the case of scattering on acoustic phonons ($`\alpha =3`$) and $`T\omega _c`$, we thus arrive at $`\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}T^7`$. This temperature dependence resembles that for the pure dephasing of an isolated state of a point center, derived by McCumber and Sturge. McCumber63 The only difference is that the exciton dephasing rate undergoes suppression by a factor of $`N+1`$ due to the motional narrowing effect. We note that this narrowing is not observed for one-phonon-assisted dephasing \[see Eqs. (25) and (26)\]. The $`T^2`$-scaling of $`\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}`$ in the high-temperature limit ($`T\omega _c`$) results from the square of the mean phonon occupation number involved in Eq. (20). To conclude this section, we stress that the $`T^{2\alpha +1}`$ and $`T^2`$ scaling relations of the pure dephasing rate with temperature obtained here, also hold for disordered aggregates, because this result is determined only by the two-vibration spectral density $`^{(2)}(\omega ,\omega ^{})`$. The suppression factor, however, will then be determined by the exciton localization size, rather than the chain length. ### IV.3 Inelastic two-phonon-assisted dephasing Finally, we analyze the dephasing rate of the superradiant state ($`\nu =1`$) resulting from the two-phonon inelastic scattering of excitons. This rate, which will be denoted as $`\mathrm{\Gamma }_{\mathrm{inel}}^{(2)}`$, is determined by the sum of scattering rates to all higher states, $`\mathrm{\Gamma }_{\mathrm{inel}}^{(2)}=(1/2)_{\mu 1}W_{\mu 1}^{(2)}`$. Using Eqs. (18) and (19), and making the same assumptions as in the case of one-phonon-assisted dephasing (Sec. IV.1), we obtain $`\mathrm{\Gamma }_{\mathrm{inel}}^{(2)}`$ $`=`$ $`{\displaystyle \frac{W_0^{(2)}}{4\pi }}\left({\displaystyle \frac{T}{J}}\right)^{2\alpha +\frac{3}{2}}{\displaystyle _0^{\omega _c/T}}dx{\displaystyle _0^{\omega _c/T}}dy`$ (29) $`\times `$ $`{\displaystyle \frac{x^\alpha y^\alpha }{(e^x1)(e^y1)}}`$ $`\times `$ $`\left[{\displaystyle \frac{2e^y}{\sqrt{xy}}}\mathrm{\Theta }(xy)+{\displaystyle \frac{1}{\sqrt{x+y}}}\right].`$ In the low-temperature limit, $`T\omega _c`$ ($`\omega _c/T\mathrm{}`$), Eq. (29) yields $$\mathrm{\Gamma }_{\mathrm{inel}}^{(2)}=\frac{\kappa W_0^{(2)}}{4\pi }\left(\frac{T}{J}\right)^{2\alpha +\frac{3}{2}},$$ (30) where the numerical factor $`\kappa `$ is given by the double integral in Eq. (29), with both upper limits replaced by infinity. Comparing Eq. (30) with Eq. (28a), we see that $`\mathrm{\Gamma }_{\mathrm{inel}}^{(2)}`$ is characterized by a steeper temperature dependence than $`\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}`$. In particular, for $`\alpha =3`$ the rate $`\mathrm{\Gamma }_{\mathrm{inel}}^{(2)}T^{15/2}`$. This means that at higher temperatures the inelastic two-phonon channel of dephasing can compete with the elastic one (Sec. V.2). In the high-temperature limit, $`T\omega _c`$ ($`\omega _c/T0`$), the rate $`\mathrm{\Gamma }_{\mathrm{inel}}^{(2)}`$ is proportional to $`T^2`$, which is the same scaling relation as in the case of pure dephasing \[Eq. (28b)\]. ## V Disordered aggregates In this section, we will analyze the temperature dependence of the various dephasing contributions in the presence of disorder and their effect on the width of the J-band. As we will see (and already anticipated in Sec. III.1), the dephasing rates are spread over a wide region, in particular at low temperature. As a consequence, it is not clear a priori what value (mean, typical, or other) of these rates should be related to the homogeneous width of the absorption spectrum. Therefore, the width of the J-band was obtained by direct simulation of the absorption spectrum, using the calculated dephasing rates to broaden each of the exciton transitions in the band. Explicitly, we have $$A(E)=\frac{1}{N}\underset{\nu }{}\frac{F_\nu }{\pi }\frac{\mathrm{\Gamma }_\nu }{(EE_\nu )^2+\mathrm{\Gamma }_\nu ^2}.$$ (31) where $`F_\nu =(_{n=1}^N\phi _{\nu n})^2`$ is the dimensionless oscillator strength of the $`\nu `$th exciton state and $`\mathrm{\Gamma }_\nu =\gamma _\nu /2+\mathrm{\Gamma }_\nu ^{(1)}+\mathrm{\Gamma }_\nu ^{(2)}`$ is the total homogeneous width of this state. Here, $`\mathrm{\Gamma }^{(1)}`$ and $`\mathrm{\Gamma }^{(2)}`$ are given by Eqs. (15) and (21), respectively, and $`\gamma _\nu =\gamma _0F_\nu `$ is the radiative decay rate of state $`\nu `$ ($`\gamma _0`$ denotes the radiative constant of a monomer). In all the simulations, we will assume the limit $`\omega _cT`$, which implies that $`\omega _c\mathrm{}`$ in Eqs. (14) and (19). As before, the angular brackets denote the average over the random realizations of the site energies $`\{\epsilon _n\}`$. The resulting J-bandwidth $`\mathrm{\Delta }`$ was determined as the full width at half maximum of the thus calculated absorption spectrum. A similar approach has been used Basko03 to simulate the absorption spectrum of THIATS aggregates Scheblykin96 at room temperature. ### V.1 One-phonon-assisted dephasing In order to study fluctuations in the dephasing rates, we analyzed their statistics, focusing on the lowest exciton state for each randomly generated disorder realization. This choice was motivated by the fact that the low-lying states dominate the absorption spectrum. The dephasing rate of the lowest state is denoted $`\mathrm{\Gamma }_{}^{(1)}`$; its distribution, collected by considering $`3\times 10^4`$ disorder realizations for chains of $`N=500`$ molecules and a disorder strength $`\sigma =0.135J`$ is presented in Figs. 1(a) and 1(b). In generating these figures, we used a one-vibration spectral density of the form Eq. (14), with $`\alpha =3`$. Furthermore, the rates were calculated for $`T=0.06J`$ \[Fig. 1(a)\] and $`T=0.23J`$ \[Fig. 1(b)\]; for the protoptypical aggregates of pseudoisocyanine ($`J=600`$ cm<sup>-1</sup>), this agrees with temperatures of 50 K and 200 K, respectively. We note that the horizontal axis of the distributions is scaled by the average value $`\overline{\mathrm{\Gamma }}_{}^{(1)}`$, so that the value of $`W_0^{(1)}`$ does not affect the figures. Of course, this scaling renders the axis temperature dependent, because $`\overline{\mathrm{\Gamma }}_{}^{(1)}`$ strongly depends on $`T`$, as we will see below (Fig. 3). These figures clearly demonstrate that the relative spread in $`\mathrm{\Gamma }_{}^{(1)}`$ may be considerable, in particular at low temperatures. This may be understood from the local band-edge level structure of the disordered tight-binding Hamiltonian. Malyshev91 For temperatures smaller than the J-bandwidth, the exciton scatters between discrete levels in the vicinity of the band edge that are localized in the same region of the chain. Both the energy spacing between these states and their localization size undergo large fluctuations: $`\delta E_{\mu \nu }E_{\mu \nu }`$ and $`\delta (_{n=1}^N\phi _{\mu n}^2\phi _{\nu n}^2)_{n=1}^N\phi _{\mu n}^2\phi _{\nu n}^2`$Malyshev01 As a consequence, $`\delta W_{\mu \nu }^{(1)}W_{\mu \nu }^{(1)}`$. When the temperature is increased, the exciton in the lowest state may scatter to many higher-lying states, which often are delocalized over an appreciable part of the chain. This smears the fluctuations that occur in the scattering rates between individual states, leading to a decrease in the relative spread of the dephasing rate of the lowest state. We next turn to the absorption band calculated according to Eq. (31), neglecting the role of two-phonon scattering \[$`W_0^{(2)}=0`$\]. In Fig. 2(a) this band is plotted for three temperatures at a fixed disorder strength of $`\sigma =0.2J`$. For the one-phonon spectral density we used the form Eq. (14) with $`W_0^{(1)}=25J`$ and $`\alpha =3`$, and took $`\gamma _0=1.5\times 10^5J`$, which is typical for J-aggregates of polymethine dyes. Chains of $`N=500`$ molecules were considered. The simulated spectra clearly demonstrate the thermal broadening, caused by growing homogeneous widths of the individual exciton transitions. At low temperature, the homogeneous broadening is negligible, the J-band is inhomogeneous, with a width that is determined by the disorder strength. With growing temperature, the J-band becomes more homogeneous, as is apparent from the fact that it gets more symmetric. In Fig. 2(b) we plotted by symbols the simulated J-band width $`\mathrm{\Delta }(T)`$ as a function of temperature for three values of the disorder strength: $`\sigma =0.1J`$, $`0.2J`$, and $`0.3J`$ \[all other parameters were taken as in Fig. 2(a)\]. As is seen, the $`\mathrm{\Delta }(T)`$ shows a plateau at the value of the inhomogeneous width, $`\mathrm{\Delta }(0)=0.04J`$, $`0.1J`$, and $`0.18J`$, respectively. Beyond these plateaus, $`\mathrm{\Delta }(T)`$ goes up quite steeply, reflecting the fact that the homogeneous (dynamic) broadening becomes dominant. To accurately extract at low temperatures the small homogeneous contribution to the total width, we generated up to $`4\times 10^5`$ disorder realizations. At higher temperatures, this number could be restricted to 4000, owing to the reduction of the relative fluctuations (cf. Fig. 1). Inspired by the analytically obtained power-laws for the one-phonon-assisted dephasing rate as a function of temperature for homogeneous aggregates (Sec. IV.1), we considered a parametrization of the form $$\mathrm{\Delta }(T)=\mathrm{\Delta }(0)+aW_0^{(1)}\left(T/J\right)^p$$ (32) for the total band width in disordered aggregates. It turned out that the calculated line widths as a function of temperature could be fitted very well by the relation (32) \[curves in Fig. 2(b)\]. The corresponding fit parameters are $`a=1.24`$ and $`p=4.16`$ for $`\sigma =0.1J`$, $`a=1.32`$ and $`p=4.29`$ for $`\sigma =0.2J`$, and $`a=1.20`$ and $`p=4.27`$ for $`\sigma =0.3J`$. The scaling relation Eq. (32) turns out to hold over an even wider range of $`\sigma `$ and $`W_0`$ values.Heijs05 This implies that, although the J-band is built up from a distribution of exciton states with different dephasing rates, the total width $`\mathrm{\Delta }(T)`$ may effectively be separated in an inhomogeneous width, $`\mathrm{\Delta }(0)`$, and a dynamic contribution. We note that the fitting exponent $`p`$ is larger than the value 3.85 found in the absence of disorder (Sec. IV.1). This increase results from downward scattering processes between optically dominant exiton states, which are possible in the presence of disorder, but not for the superradiant state in the homogeneous chain. This claim may be substantiated by considering the dephasing rates of the lowest exciton state of each disorder realization. We numerically generated the average of this quantity, $`\overline{\mathrm{\Gamma }}_{}^{(1)}`$, from $`3\times 10^4`$ disorder realizations for chains of $`N=500`$ molecules with $`\sigma =0.1J`$, $`W_0^{(1)}=25J`$, and $`\alpha =3`$. This average is shown as a function of temperature in Fig. 3 (diamonds), together with the dynamic contribution to the total J-bandwidth, $`\mathrm{\Delta }(T)\mathrm{\Delta }(0)`$ (solid line), and the dephasing rate $`\mathrm{\Gamma }_1^{(1)}`$ of the superradiant state for a homogeneous chain of the same length (squares). As is seen, the dynamic part $`\mathrm{\Delta }(T)\mathrm{\Delta }(0)`$ has a larger exponent $`p=4.16`$ than $`\overline{\mathrm{\Gamma }}_{}^{(1)}`$ $`(p=3.85)`$. It is remarkable, however, that $`\overline{\mathrm{\Gamma }}_{}^{(1)}`$ and $`\mathrm{\Gamma }_1^{(1)}`$ display almost identical behavior, at least in the relevant region $`T\mathrm{\Delta }(0)`$, where the homogeneous contribution to the J-bandwidth is noticeable. For $`T\mathrm{\Delta }(0)`$, where $`\mathrm{\Gamma }_{}^{(1)}`$ undergoes significant fluctuations (see Fig. 1), $`\overline{\mathrm{\Gamma }}_{}^{(1)}`$ turns out to be much smaller than $`\mathrm{\Gamma }_1^{(1)}`$. In the high-temperature regime, these fluctuations are washed out and the two quantities become almost identical. So far, we have only presented numerical results for a one-phonon spectral density Eq. (14) with the power $`\alpha =3`$, which corresponds to a Debye model for the host vibrations. To end this subsection, we will address the effect of changing the spectral density. First, we consider the effect of the value for $`\alpha `$. The diamonds in Fig. 4 present our results for the temperature dependence of the width $`\mathrm{\Delta }(T)`$ obtained for $`\alpha =1`$ and $`W_0=5J`$. As before, we found that these data may be fitted by a simple power-law of the form Eq. (32) (dashed line). In this case we find $`p=1.9`$, which, again, is slightly larger than the value 3/2 found from Eq. (25), due to the correction of the dispersion relation arising from the long-range dipole-dipole interactions and downward scattering processes that contribute to the total J-bandwidth. Clearly, these data demonstrate the sensitivity of the temperature dependence of the total linewidth to the power $`\alpha `$. Interestingly, it turns out that while the temperature dependence of the J-bandwidth is sensitive to the overall frequency scaling of the spectral density, it is not sensitive to fluctuations of this scaling around an average power-law. To demonstrate this, we have considered a spectral density of the form $`^{(1)}(\omega )=W_0^{(1)}(\omega /J)^3[1+\mathrm{sin}(2\pi \omega /\stackrel{~}{\omega })]`$, which only on average exhibits an $`\omega ^3`$-dependence. The results for the width $`\mathrm{\Delta }(T)`$ obtained for $`W_0^{(1)}=25J`$ and $`\stackrel{~}{\omega }=J/6`$ are presented as squares in Fig. 4. Remarkably, these results are indistinguishable from those obtained without fluctuations (i.e., $`\stackrel{~}{\omega }=\mathrm{}`$: triangles). The reason is that at elevated temperatures the function $`\omega ^3\overline{n}(\omega )`$ varies slowly on the scale of $`\stackrel{~}{\omega }`$, so that the modulating function $`1+\mathrm{sin}(2\pi \omega /\stackrel{~}{\omega })`$ may be replaced by its average value, which equals unity. ### V.2 Two-phonon-assisted dephasing As we have seen in Sec.III.2, the two-phonon-assisted dephasing rate $`\mathrm{\Gamma }_\nu ^{(2)}`$ consists of four contributions, one of which is elastic (indicated as “pure”, as it is responsible for pure dephasing), while the other three are inelastic and are indicated as downward ($``$), cross ($``$), and upward ($``$). In Figs. 1(c)- 1(f), 5, and 6 we present results for the statistics of these various contributions for disordered aggregates. In all cases, we used chains of 500 molecules, a disorder strength of $`\sigma =0.135J`$, and a two-vibration spectral density of the form Eq. (19), with $`\alpha =3`$. The statistics are presented for the lowest exciton state in each one of $`3\times 10^4`$ randomly generated disorder realizations. For this state the downward contribution vanishes and the two-phonon-assisted dephasing rate reads $`\mathrm{\Gamma }_\nu ^{(2)}=\mathrm{\Gamma }_{}^{(2)}+\mathrm{\Gamma }_{}^{(2)}+\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}`$. These three remaining contributions were calculated using the expressions derived in the Appendix. From Figs. 1(c) and 1(d) we observe that at a given temperature, the relative spread in $`\mathrm{\Gamma }_{}^{(2)}`$ is much smaller than that in $`\mathrm{\Gamma }_{}^{(2)}`$. The reason is that in a two-phonon-assisted upward process the exciton in the lowest state scatters to more higher-energy states than in a cross process. As a result, fluctuations in $`\mathrm{\Gamma }_{}^{(2)}`$ are suppressed more than those in $`\mathrm{\Gamma }_{}^{(2)}`$. Upon heating, the spread of both $`\mathrm{\Gamma }_{}^{(2)}`$ and $`\mathrm{\Gamma }_{}^{(2)}`$ reduces, which has the same explanation as given for this effect in the case of one-phonon-assisted dephasing (Sec. V.1). The distribution of $`\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}`$ does not depend on temperature at all, because, according to Eq. (20), the rate $`\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}`$ fluctuates exclusively due to fluctuations in the inverse participation ratio $`_{n=1}^N\phi _{\nu n}^4`$. As the latter quantity is subject to large fluctuations, Malyshev01 this also explains the large relative spread in $`\mathrm{\Gamma }_{\mathrm{pure}}^{(2)}`$. In Fig. 6 we plotted (on a log-log scale) the temperature dependence of the mean values $`\overline{\mathrm{\Gamma }}_{}^{(2)}`$, $`\overline{\mathrm{\Gamma }}_{}^{(2)}`$, and $`\overline{\mathrm{\Gamma }}_{\mathrm{pure}}^{(2)}`$, obtained from averaging over these rates for the lowest exciton states in the simulations discussed above. This figure nicely shows the relative importance of the different dephasing channels. We clearly see that $`\overline{\mathrm{\Gamma }}_{}^{(2)}\overline{\mathrm{\Gamma }}_{}^{(2)},\overline{\mathrm{\Gamma }}_{\mathrm{pure}}^{(2)}`$, i.e., the inelastic channel of dephasing due to double phonon absorption is inefficient, at all temperatures. More importantly, we observe that the cross channel of inelastic two-phonon-assited dephasing successfully competes with the pure dephasing contribution. At low temperatures, pure dephasing dominates, while at higher temperatures the inelastic cross process is more important. This correlates well with our findings for disorder-free aggregates \[see discussion below Eq. (30)\]. For the disorder strength considered here, the cross-over occurs at $`T_0=0.12J`$ ($`100`$ K for $`J=600`$ cm<sup>-1</sup>). We note that when considering two-phonon scattering, one often restricts to modelling the elastic process (see, e.g., Ref. Kuhn97, ). From the above we see that this is not justified at elevated temperatures. Despite the fact that not all three curves in Fig. 6 are exactly straight lines (deviations occur in the unimportant low-temperature part), they all can be fitted very well by a power-law, $`a(T/J)^p`$. In doing so, we obtained $$\overline{\mathrm{\Gamma }}_{}^{(2)}=0.18W_0^{(2)}\left(T/J\right)^{8.3},$$ (33a) $$\overline{\mathrm{\Gamma }}_{}^{(2)}=9.01W_0^{(2)}\left(T/J\right)^{7.9},$$ (33b) $$\overline{\mathrm{\Gamma }}_{\mathrm{pure}}^{(2)}=1.88W_0^{(2)}\left(T/J\right)^7.$$ (33c) We recall that the exponent $`p=7`$ in the last formula is an exact result for $`\alpha =3`$ and $`\omega _cT`$, as was argued at the end of Sec. IV.2 already. As in the case of one-phonon scattering, we see that the inelastic two-phonon dephasing rates exhibit a steeper temperature dependence than for the homogeneous aggregate with nearest-neighbor interactions \[Eq. (30)\]. ## VI Comparison to experiment and discussion It is of interest to see to what extent the model we presented here is able to explain the temperature dependence of the homogeneous broadening in molecular aggregates. As mentioned in the Introduction, Renge and Wild Renge97 found that the total J-bandwidth $`\mathrm{\Delta }(T)`$ of PIC-Cl and PIC-F over a wide temperature range (from 10 K to 300 K) follows a power-law scaling as in Eq. (32). Although the power reported by these authors ($`p=3.4`$) is smaller than the ones we derived in Sec. V.1, from direct comparison to the experimental data we have found that a model of one-phonon scattering with a spectral density given by Eq.(14) with $`\alpha =3`$ and $`\omega _c\mathrm{}`$, yields an excellent quantitative explanation of the experiments over the entire temperature range, both for the shape and the width of the J-band. The same turns out to be true for the J-bandwidth of PIC-Br, measured between 1.5 K and 180 K. Fidder90 In all these fits, $`\sigma `$ and $`W_0`$ were the only two free parameters that could be adjusted to optimize comparison to experiment. Details will be published elsewhere,Heijs05 together with a fit of the much debatedFidder90 ; deBoer89 ; Spano90 ; Fidder95 ; Potma98 temperature dependence of the fluorescence lifetime of these aggregates. Here, we present an explicit comparison to the hole-burning data reported by Hirschmann and Friedrich.Hirschmann89 Using this technique, they measured the homogeneous width of the exciton states in the center of the J-band for PIC-I over the temperature range 350 mK to 80 K. Their data for the holewidth $`\mathrm{\Gamma }`$ are reproduced as triangles in Fig. 7. The solid line shows our fit to these data, obtained by simulating disordered chains of $`N=250`$ molecules with a one-phonon spectral density of the form Eq. (14) with $`\alpha =3`$ and $`\omega _c\mathrm{}`$. The resonant interaction strength and the monomer radiative rate were chosen at the accepted values of $`J=600`$ cm<sup>-1</sup> and $`\gamma _0=1.5\times 10^5J=2.7\times 10^8`$ s<sup>-1</sup>, respectively. Thus, the only free parameters were the disorder strength $`\sigma `$ and scattering strength $`W_0^{(1)}`$. First, we fixed the value of $`\sigma `$ by fitting the low-temperature (4 K) absorption spectrum, where the homogeneous broadening may be neglected. This yielded $`\sigma =0.21J`$. Next, $`W_0^{(1)}`$ was adjusted such that the measured growth of the hole width was reproduced in an optimal way. Thus, we found $`W_0^{(1)}=180J`$. For each given temperature the simulated hole width was obtained as the average of the dephasing rate of the excitons found in an interval of 0.06 cm<sup>-1</sup> around the center of the simulated J-band at that temperature. Taking into account the fair amount of scatter in the experimental data, we conclude from Fig. 7 that our model yields a good fit to the measurements. As mentioned above, the same holds for the J-bandwidth in aggregates of PIC-Cl, PIC-F, and PIC-Br. This yields valuable information about the dominant mechanism of dephasing in these materials. We conclude that this mechanism is one-phonon-assisted scattering of excitons on vibrations in the host characterized by a spectral density which (on average) scales as $`\omega ^3`$; this scaling yields a natural explanation of the power-law thermal broadening of the J-band found in various experiments. We stress that it is impossible to fit the experimental data with a spectral density that is constant or scales linearly with $`\omega `$, as that yields considerably different power laws for the width \[cf. Fig. 4\]. The $`\omega ^3`$ scaling of the spectral density needed to fit the experiments strongly suggests that acoustic phonons dominate the scattering process. Thus, the spectral width measured over a broad temperature range is an excellent probe for the scattering mechanism. It is appropriate to comment on the value of $`W_0^{(1)}`$ obtained from our fit, which seems to be very large. It should be kept in mind that $`W_0^{(1)}`$ is a phenomenological scattering strength, which combines several microscopic material properties \[see discussion below Eq. (14)\]. Most importantly, the value found here is consistent with a perturbative treatment of the scattering process: the scattering rates between the optically dominant states obtained from it turned out to be much smaller than their energy separation. We finally address an alternative mechanism of dephasing, namely scattering on local vibrations belonging to the aggregate. This has been suggested by several authors based on activation-law fits of the measured homogeneous contribution to the J-bandwidth. Fidder90 ; Hirschmann89 To get an estimate whether this is a reasonable mechanism, let us neglect the disorder and make the nearest-neighbor approximation for the resonance interactions $`J_{nm}`$. Then, $`E_\nu `$ and $`\phi _{\nu n}`$ are given by Eq. (22). Close to the lower exciton band edge, the region of our interest, $`E_\nu =2J+J\pi ^2\nu ^2/(N+1)^2`$. Furthermore, let us parameterize the spectral function of a local vibration of frequency $`\omega _0`$ as $`(\omega )=2\pi V_0^2\delta (\omega \omega _0)`$, where $`V_0`$ is the coupling constant to the excitons. We are interested in the dephasing rate $`\mathrm{\Gamma }_1`$ of the superradiant state $`\|\nu =1`$. Using the above simpifications and replacing in Eq. (15) the summation over exciton states by an integration, one easily arrives at $$\mathrm{\Gamma }_1=V_0^2\frac{\overline{n}(\omega _0)}{(J\omega _0)^{1/2}}.$$ (34) Comparing this result to the activation law $`b\mathrm{exp}(\omega _0/T)`$, used in Refs. Fidder90, and Hirschmann89, to fit the experimental data, one obtains $`V_0^2=b(J\omega _0)^{1/2}`$. Substituting the values $`b=3000`$ cm<sup>-1</sup> and $`\omega _0=330`$ cm<sup>-1</sup> from Ref. Hirschmann89, , and $`J=600`$ cm<sup>-1</sup>, we obtain as estimate for the exciton-phonon coupling $`V_01100`$ cm<sup>-1</sup>. This is an enormously large number. In particular, the Stokes losses $`S=V_0^2/\omega _03700`$ cm<sup>-1</sup> turn out to be much larger than the resonant interaction $`J=600`$ cm<sup>-1</sup>. Under these conditions, a strong exciton self-trapping is to be expected, resulting in a reduction of the exciton bandwidth by a factor of $`\mathrm{exp}(S/\omega _0)1`$Rashba82 Coherent motion of the exciton, even over a few molecules, is then hardly possible: any disorder will destroy it. This observation is not consistent with the widely accepted excitonic nature of the aggregate excited states, as corroborated by many optical and transport measurements.Kobayashi96 ; Knoester02 ## VII Concluding remarks In this paper we presented a theoretical study of the temperature dependence of the exciton dephasing rate in linear J-aggregates and the resulting width of the total absorption band (the J-band). As dephasing mechanism we considered scattering of the excitons on vibrations of the host matrix, taking into account both one- and two-vibration scattering. The excitons were obtained from numerical diagonalization of a Frenkel exciton Hamiltonian with energy disorder and their dephasing rates were subsequently calculated using a perturbative treatment of the exciton-vibration interaction (Fermi Golden Rule). In the absence of disorder, the lowest (superradiant) exciton state dominates the absorption spectrum. As a result, the dephasing rate of this state directly gives the homogeneous width of the J-band. We analytically calculated the temperature dependence of this homogeneous width for both one- and two-vibration scattering, assuming a Debye-like model for the host vibration density of states. It turned out that in all cases the homogeneous width obeys a power-law as a function of temperature, with the value of the exponent depending on the shape of the low-energy part of the vibronic spectrum (Sec. IV). In the presence of disorder the optically dominant exciton states still reside close to the bottom of the band, but their energies are now spread and their wave functions become localized on finite segments of the chain. We have found that the various one- and two-phonon-induced contributions to the dephasing rate undergo significant fluctuations, because the exciton energies and overlap integrals vary considerably from one disorder realization to the other (Sec. V). As a consequence, one cannot use the dephasing rates of individual states to characterize the homogeneous broadening of the J-band. Instead, we simulated the total J-band and showed that its width effectively separates in an inhomogeneous (zero-temperature) contribution and a dynamic (homogeneous) part. The latter scales with temperature according to a power law \[Eq. (32)\], with an exponent that is somewhat larger than the one found for the homogeneous broadening in the absence of disorder. We also showed that amongst the two-vibration scattering processes, inelastic channels will at elevated temperatures dominate the usually considered pure-dephasing contribution. Finally, from comparison to absorption and hole-burning experiments (Sec. VI), we found that the dominant mechanism of dephasing for J-aggregates lies in one-phonon scattering of excitons on vibrations of the host matrix, characterized by a spectral density which (on average) scales like the third power of the phonon frequency. This suggests that acoustic phonons of the host play an important role in the scattering process. All temperature dependent data available to date, are consistent with this picture. By contrast, we argued that the previously suggested mechanism of scattering on local vibrations of the aggregate, leading to an activated thermal behavior, is not consistent with the overwhelming amount of evidence that the optical excitations in J-aggregates have an excitonic character. ## Appendix A Two-phonon scattering rates In this Appendix we present expressions for the two-phonon-assisted scattering rate $`W_{\mu \nu }^{(2)}`$, starting from Eq. (18). After substituting the two-vibration spectral density given in Eq. (19) and performing several algebraic manipulations, we arrive at $`W_{\mu \nu }^{(2)}`$ $`=`$ $`W_0^{(2)}\left({\displaystyle \frac{T}{J}}\right)^{2p+1}{\displaystyle \underset{n=1}{\overset{N}{}}}\phi _{\mu n}^2\phi _{\nu n}^2`$ $`\times `$ $`\left[F_{}^{(2)}(\omega _{\mu \nu })+2F_{}^{(2)}(\omega _{\mu \nu })+F_{}^{(2)}(\omega _{\mu \nu })\right],`$ where the temperature dependent functions $`F_{}^{(2)}(\omega _{\mu \nu })`$, $`F_{}^{(2)}(\omega _{\mu \nu })`$, and $`F_{}^{(2)}(\omega _{\mu \nu })`$ distinguish between the scattering processes in which two phonons are emitted ($``$), one is absorbed and one is emitted ($``$), and two phonons are absorbed ($``$). They are given by $`F_{}^{(2)}(\omega _{\mu \nu })`$ $`=`$ $`{\displaystyle _0^{\omega _c/T}}𝑑x{\displaystyle _0^{\omega _c/T}}𝑑y`$ (36a) $`\times `$ $`f_{}(x,y)\delta \left({\displaystyle \frac{\omega _{\mu \nu }}{T}}+x+y\right),`$ $`F_{}^{(2)}(\omega _{\mu \nu })`$ $`=`$ $`{\displaystyle _0^{\omega _c/T}}𝑑x{\displaystyle _0^{\omega _c/T}}𝑑y`$ (36b) $`\times `$ $`f_{}(x,y)\delta \left({\displaystyle \frac{\omega _{\mu \nu }}{T}}x+y\right),`$ $`F_{}^{(2)}(\omega _{\mu \nu })`$ $`=`$ $`{\displaystyle _0^{\omega _c/T}}𝑑x{\displaystyle _0^{\omega _c/T}}𝑑y`$ (36c) $`\times `$ $`f_{}(x,y)\delta \left({\displaystyle \frac{\omega _{\mu \nu }}{T}}xy\right),`$ where, after changing to dimensionless integration variables $`x=\omega _q/T`$ and $`y=\omega _q^{}/T`$, we introduced the auxiliary functions $`f_{}(x,y)`$ $`=`$ $`{\displaystyle \frac{x^pe^x}{e^x1}}{\displaystyle \frac{y^pe^y}{e^y1}},`$ (37a) $`f_{}(x,y)`$ $`=`$ $`{\displaystyle \frac{x^p}{e^x1}}{\displaystyle \frac{y^pe^y}{e^y1}},`$ (37b) $`f_{}(x,y)`$ $`=`$ $`{\displaystyle \frac{x^p}{e^x1}}{\displaystyle \frac{y^p}{e^y1}}.`$ (37c) We note that for $`\mu =\nu `$ ($`\omega _{\mu \nu }=0`$), the only non vanishing term is $`F_{}^{(2)}(0)`$, which describes the elastic channel of scattering (pure dephasing). We now further analyze the three contributions to the scattering rate. As $`_{}^{(2)}(\omega _{\mu \nu })`$ describes the emission of two vibrational quanta, we have $`\omega _{\mu \nu }<0`$. Thus, performing the $`y`$ integration, we obtain $`F_{}^{(2)}(\omega _{\mu \nu })={\displaystyle _0^{\omega _{\mu \nu }/T}}dxf_{}(x,{\displaystyle \frac{\omega _{\mu \nu }}{T}}x),`$ (38a) if $`\omega _{\mu \nu }<\omega _c`$, and $`F_{}^{(2)}(\omega _{\mu \nu })={\displaystyle _{(\omega _{\mu \nu }+\omega _c)/T}^{\omega _c/T}}dxf_{}(x,{\displaystyle \frac{\omega _{\mu \nu }}{T}}x),`$ if $`\omega _c<\omega _{\mu \nu }<2\omega _c`$, while $`F_{}^{(2)}(\omega _{\mu \nu })=0`$ otherwise. Next, the contribution which involves the emission and absorption of one vibrational quantum, is given by $`F_{}^{(2)}(\omega _{\mu \nu })={\displaystyle _{\omega _{\mu \nu }/T}^{\omega _c/T}}dxf_{}(x,{\displaystyle \frac{\omega _{\mu \nu }}{T}}+x),`$ (39a) if $`0\omega _{\mu \nu }<\omega _c`$. If $`0<\omega _{\mu \nu }<\omega _c`$, $`F_{}^{(2)}(\omega _{\mu \nu })={\displaystyle _0^{(\omega _c+\omega _{\mu \nu })/T}}dxf_{}(x,{\displaystyle \frac{\omega _{\mu \nu }}{T}}+x),`$ (39b) and $`F_{}^{(2)}(\omega _{\mu \nu })=0`$ otherwise. Finally, the contribution that results from the absorption of two vibrational quanta ($`\omega _{\mu \nu }>0`$), yields $`F_{}^{(2)}(\omega _{\mu \nu })={\displaystyle _0^{\omega _{\mu \nu }/T}}dxf_{}(x,{\displaystyle \frac{\omega _{\mu \nu }}{T}}x),`$ (40a) if $`0<\omega _{\mu \nu }<\omega _c`$, and $`F_{}^{(2)}(\omega _{\mu \nu })={\displaystyle _{(\omega _{\mu \nu }\omega _c)/T}^{\omega _c/T}}dxf_{}(x,{\displaystyle \frac{\omega _{\mu \nu }}{T}}x),`$ (40b) if $`\omega _c<\omega _{\mu \nu }<2\omega _c`$, while $`F_{}^{(2)}(\omega _{\mu \nu })=0`$ otherwise.
warning/0506/hep-ph0506138.html	ar5iv	text	# 1 Introduction ## 1 Introduction History of exotic hadrons is as old as that of the quark model , and the subject has been studied for long time . Yet the recent observation of the evidence of the pentaquark particle $`\mathrm{\Theta }^+`$ has triggered enormous amount of research activities both in experimental and theoretical hadron physics . Baryons containing five valence quarks are totally new form of hadrons. The importance of knowing the nature of multi-quark states lies, for instance, in understanding the origin of matter. It is believed that in the early stage of the universe, matter was highly dense forming the quark matter. A natural question would then be what is the mechanism of the transition from that to the present hadronic world consisting of ordinary mesons and baryons. At this moment, the existence of the pentaquarks is still the most important issue. The whole discussions below are, therefore, based on this assumption. From the hadron physics point of view, the understanding of five quark systems, if they exist as (quasi-)stable states, will give us more information on the dynamics of non-perturbative QCD, such as confinement of colors and chiral symmetry breaking. Many ideas have been proposed attempting to explain the unique features of $`\mathrm{\Theta }^+`$. As it has turned out and will be discussed in this note, however, the current theoretical situation is not yet settled at all, having revealed that our understanding of hadron physics would be much poorer than we have thought . We definitely need more solid ideas and methods to answer the related questions. Turning to the specific interest in $`\mathrm{\Theta }^+`$, its would-be light mass and narrow width are the issues to be understood, together with the determination of its spin and parity. In particular, the information of parity is important, since it reflects the internal motion of the constituents. In this lecture the following materials are discussed, with emphasis on theoretical methods. * A quick overview on experimental situation (section 2). * Some basics of theoretical models; the quark model, chiral soliton model , and somewhat general treatment based on the flavor SU(3) symmetry as well as a brief view over lattice and QCD sum rule studies (section 3). * Decay of $`\mathrm{\Theta }^+`$ (section 4) * Production reactions including photo and hadronic productions (section 5). Through these discussions, we consider a possibility of $`\mathrm{\Theta }^+`$ with $`J^P=3/2^{}`$ as one of likely candidates for a pentaquark state. There are many interesting topics which can not be discussed in this note. For readers who are interested in more details, please refer to the proceedings of the workshop PENTAQUARK04 and references in there . ## 2 Experiments The first observation was made by the LEPS group at SPring-8 lead by T. Nakano . The backward compton-scattered photon of energy 2.4 GeV produced at SPring-8 was used to hit a neutron target inside a carbon nucleus to produce a strangeness and antistrangeness pair ($`K^+`$ and $`K^{}`$). The Fermi motion corrections were carefully analyzed, and then a missing mass analysis was performed for the $`K^+n`$ final state. They have seen an excess in the $`K^+n`$ invariant mass spectrum over the background at 4.6$`\sigma `$ level in the energy region 1.54 GeV. The width of the peak was as narrow as or less than the experimental resolution ($`25`$ MeV). The peak was then identified with the exotic pentaquark state of strangeness $`S=+1`$. The absence of the similar peak structure in the $`K^+p`$ system suggests the isospin of the state is likely to be $`I=0`$. The spectrum of the LEPS experiment is shown in Fig. 1, where the peak around 1.54 GeV is the first signal of the exotic particle . Immediately after the announce of this results, many positive signals follow . Among them, the existence of another exotic baryon of strangeness $`1`$, $`\mathrm{\Xi }^{0,,}`$, was also reported . Major results of experiments so far are summarized in Table 1. From there, one can recognize that there is fluctuation in absolute values of the mass of $`\mathrm{\Theta }^+`$, from 1520 to 1550 MeV. It is often said that the fluctuation of order 30 MeV is large; it is about 30 % level if measured from the $`KN`$ threshold. After many positive signals were reported, negative results followed also, mostly from the analysis of high energy experiments . These are also summarized in Table 1. At this moment there is not a theory consistently explain these data. The high energy experiments have much higher statistics than the low energy experiments and should be taken seriously. If $`\mathrm{\Theta }^+`$ exists and can be seen only in the low energy (mostly in photoproductions) experiments, one needs to understand the production mechanism . A model for the suppression at high energies was proposed by Titov et al. . If it does not, we also need to understand what the signals in the low energy experiments are for. Very recently, CLAS (g11) reported the null result in the reaction $`\gamma p\overline{K}^0K^+n`$ . This has much larger statistics than the previous experiment performed at SAPHIA by about factor twenty . They extracted an upper limit of the $`\mathrm{\Theta }^+`$ production cross section, $`\sigma \underset{}{<}`$ 1 – 4 nb. The results, however, does not immediately lead to the absence of $`\mathrm{\Theta }^+`$, since there could be a large asymmetry between the reactions from the proton and neutron . In general, photoproductions are large for charge-exchange reactions, but the reaction $`\gamma p\overline{K}^0K^n`$ is not the like. Experimental studies from the neutron with higher statistics is therefore very important. ## 3 Theoretical methods In this section, we discuss the structure of the pentaquarks, especially of $`\mathrm{\Theta }^+`$. Although the pioneering work of Diakonov et. al. was performed in the chiral soliton model , it is always instructive and intuitively understandable to work in the quark model . After a brief look at the basics of pentaquark structure in the quark model, we discuss some essences of the chiral soliton model. Results of the chiral solitons are then interpreted in terms of a quark model with chiral symmetry (the chiral bag model) . After the introduction of the two models, we discuss a model independent method based on flavor SU(3) symmetry, where possible spin and parity of $`\mathrm{\Theta }^+`$ are investigated . In the last two subsections, we briefly look at the lattice QCD and QCD sum rule. ### 3.1 Constituent quark model This model has been successfully applied to the description of the conventional mesons and baryons for their masses and various transition amplitudes . In this model, a confining potential for quarks is introduced, which is usually taken to be a harmonic oscillator one, to prepare basis states as single particle states which valence quarks occupy. Then quark-quark interactions such as the spin-color interaction of one gluon exchange and the spin-flavor one of one-meson (the Nambu-Goldstone boson) exchange are introduced as residual interactions. The interaction hamiltonian is then treated either perturbatively or diagonalized within a given model space. The role of various interactions for $`\mathrm{\Theta }^+`$ has been investigated in the literatures . Here, to make discussions simple, we consider what the structure of the five-quark states are like in a confining potential. The single particle states of the harmonic oscillator potential are denoted by the principal and angular momentum quantum numbers $`(n,l)`$. Using the spectroscopic notation, we express them as $`0s,0p,1s`$ and so on. Due to the many degrees of freedom of color (3), flavor (3) and spin (2), five quarks including one anti-quark ($`\overline{s}`$) can occupy the lowest ground state simultaneously. Thus, we denote the ground state of the five quarks as $`(0s)^5`$. If one quark is excited to a $`p`$-orbit, $`(0s)^40p`$ and so on. The parity of the ground state is negative, since the antiquark carries negative parity, while the parity of the first excited state is positive. Now, let us consider flavor structure. Under the assumption of SU(3) symmetry, we need to perform irreducible decomposition of the five quark states, the direct product of four fundamental and one conjugate representations, $`3333\overline{3}=1810\overline{10}2735,`$ (1) where multiplicities are ignored on the right hand side. Among these representations on the right hand side, the isosinglet state of $`S=+1`$ appears only in the antidecuplet representation $`\overline{10}`$, which is a candidate for the SU(3) multiplet for the pentaquarks. The weight diagram of the $`\overline{10}`$ representation is shown in Fig. 2, where locations of various states are also indicated. Flavor wave functions of the antidecuplet states are easily constructed, if one notices that one of the five quark is $`\overline{3}`$. Form two $`\overline{3}`$’s out of two quark pairs (diquarks) in an antisymmetric combination, $`\overline{Q}_i=ϵ_{ijk}q_jq_k;\overline{U}[ds],\overline{D}[su],\overline{S}[ud].`$ (2) Then we can make symmetric products in terms of two diquarks and one antiquark, which generate antidecuplet members: $`\mathrm{\Theta }^+=\overline{S}\overline{S}\overline{s},N_{\overline{10}}={\displaystyle \frac{1}{\sqrt{3}}}(\overline{S}\overline{S}\overline{u}+\overline{S}\overline{U}\overline{s}+\overline{U}\overline{S}\overline{s}),`$ $`\mathrm{\Sigma }_{\overline{10}}^{}={\displaystyle \frac{1}{\sqrt{3}}}(\overline{S}\overline{U}\overline{u}+\overline{U}\overline{S}\overline{u}+\overline{U}\overline{U}\overline{s}),\mathrm{\Xi }_{\overline{10}}^{}=\overline{U}\overline{U}\overline{u}.`$ (3) There are analogous to the decuplet wave functions for ($`\mathrm{\Delta },\mathrm{\Sigma }^{},\mathrm{\Xi }^{},\mathrm{\Omega }`$). What is interesting is the average number of strange (and anti-strange) quarks in the wave functions. One can easily verify that it is 1 for $`\mathrm{\Theta }^+`$, 4/3 for $`N_{\overline{10}}`$, 5/3 for $`\mathrm{\Sigma }_{\overline{10}}`$ and 2 for $`\mathrm{\Xi }_{\overline{10}}`$. Namely, the strange quark content increases by equal amount 1/3 as the hypercharge decreases. This is a general consequence valid to the symmetric representation of SU(3). The simple counting implies that if the $`\lambda _8`$ is the only source of the SU(3) breaking as $`m_u=m_d<<m_s`$, where $`m_i`$ are constituent quark masses, the equal mass splitting of the antidecuplet baryons is expected to be $`(1/3)(m_sm_u)(1/3)\mathrm{\Delta }`$. Hence we also expect $`M(\mathrm{\Xi }_{\overline{10}})M(\mathrm{\Theta }^+)m_sm_u200\mathrm{MeV}.`$ (4) This pattern of mass splitting is shown on the left column of Fig. 3. The amount of the total mass difference $`M(\mathrm{\Xi }_{\overline{10}})M(\mathrm{\Theta }^+)`$ as shown there is significantly smaller than the one originally estimated by Diakonov et al. , the spectrum of which is shown on the right side of Fig. 3. As pointed out by Jaffe and Wilczek , the antidecuplet nucleon and sigma states mix with the corresponding octet members. We will consider the mixing effects in more detail in subsection 3.4. Here to illustrate this mixing effect in a simple case, we also show in Fig. 3 the mass pattern of the octet and the ideally mixed members. Since the ideally mixed states are classified by the strange quarks, the mass splitting between neighbors is $`\mathrm{\Delta }`$. The original prediction of the chiral soliton model is close to this in values. As seen from the figure, there is significant difference in the mass patterns in the pentaquark baryons depending on the realization of the flavor SU(3) symmetry (breaking). In general, the constituent quark model can not predict absolute values of masses. Nevertheless, if we estimate them by using typical values of constituent masses, $`m_u,m_d300`$ MeV and $`m_s500`$ MeV, we find $`M_{\mathrm{\Theta }^+}1.7`$ GeV, and other masses in accordance with the equi-distant rule. The mass of $`\mathrm{\Theta }^+`$ is larger than the observed values. In Fig. 3, however, the mass of $`\mathrm{\Theta }^+`$ is normalized. ### 3.2 Chiral solitons This model is based on the idea of Skyrme that baryons are made from weakly interacting mesons, solitons . A microscopic basis of this model is the $`1/N_c`$-expansion of QCD . The fact that there are two light flavors is also important; SU(2) isospin symmetry leads to a strongly correlated pseudoscalar pion fields under rotations in the coordinate space and isospin space. The pion field is conveniently parametrized by an SU(2) matrix as $`U(\stackrel{}{x})=\mathrm{exp}(i\stackrel{}{\tau }\stackrel{}{\pi }/f_\pi ),`$ (5) where $`f_\pi =93`$ MeV is the pion decay constant. Under the strong correlation, the hedgehog configuration is realized as a static solution where the pion field points to the radial direction, $`\stackrel{}{\pi }/f_\pi =\widehat{r}F(r)`$, where the spherical profile function $`F(r)`$ is determined by solving the classical field equation of motion. Nontrivial solutions for $`F(r)`$ define the ground states of the system. Due to nontrivial topology in the theory, different solutions $`F(r)`$ exist as classified by the winding number which is physically identified with the baryon number. The system of one baryon number describes the single nucleon sector. The hedgehog solution is a classical configuration and does not correspond to a physical nucleon state. To make a link between them, we introduce the collective variables for isospin rotations $`A(t)SU(2)`$, $`U(t,\stackrel{}{x})=A(t)U_H(\stackrel{}{x})A(t)^{},U_H(\stackrel{}{x})=\mathrm{exp}(i\stackrel{}{\tau }\widehat{r}F(r)).`$ (6) To be more precise, one needs to introduce another rotation in coordinate space, $`x_iR_{ij}x_j`$. The spatial rotation $`R`$, however, is equivalent to the isospin rotation $`A`$ due to the symmetry of the hedgehog configuration; $`A`$ and $`R`$ can not be independent degrees of freedom upon quantization. Consequently, the quantization of the $`A`$ variable leads to the wave functions which are the SU(2) $`D`$-functions for free motion in the SU(2) manifold, $`D_{t,m}^I(A)`$. Constraints then follow in the quantized states; spin and isospin must take the same values; $`J=I`$ and $`t=I_z,m=J_z`$ . When this method is applied to flavor SU(3), one finds several interesting consequences. We will state some of them without proof. The SU(3) baryonic states are written in terms of the SU(3) $`D`$-functions $`D_{YII_3;Y^RI^RI_3^R}^{(p,q)}(\alpha _1,\mathrm{}\alpha _8)`$, where the upper and lower indices label the SU(3) states, and $`\alpha _1,\mathrm{}\alpha _8`$ are the Euler angles for SU(3) rotations. A crucial observation here is that under the hedgehog ansatz, the right quantum numbers $`Y^RI^RI_3^R`$ are related to the spin and hypercharge quantum numbers. Furthermore, the Wess-Zumino term puts further constraints on the right hyper charge and the baryon number, $`(I^R,I_3^R)=(J,J_3),Y^R=B,`$ (7) where the second equation holds when $`N_c=3`$. From these, it follows that the number of states of $`Y=1`$ is $`2J+1`$. For $`\overline{10}`$, the state of $`Y=1`$ is the nucleon and so $`2J+1=2`$, or $`J=1/2`$. The parity of this state is the same as that of the nucleon, and hence the spin and parity of $`\mathrm{\Theta }^+`$ and its partner are $`J^P=1/2^+`$. The mass splitting among the multiplet $`\overline{10}`$ is once again equi-distant. In the original paper by Diakonov et al. , they determined parameters in the mass formula (see Eq. (11) below) from information of non-exotic sectors, making then prediction for the exotic baryons. The relatively low mass of $`\mathrm{\Theta }^+`$ was predicted this way prior to the observation. In their original work, the nucleon resonance $`N(1710)`$ was identified with a member of $`\overline{10}`$, to determine the equi-distant parameter $`\mathrm{\Delta }`$ 180 MeV. This pattern of the mass splitting is shown on the right side of Fig. 3. ### 3.3 Role of chiral symmetry It is instructive to make an interpretation of the results of the chiral soliton model, especially the fact that the $`1/2^+`$ state appears as the lowest state of $`\mathrm{\Theta }^+`$. As it turns out, the role of chiral symmetry is important. As we have remarked in the previous subsection, a positive parity $`\mathrm{\Theta }^+`$ requires an orbital excitation of a quark to an odd parity orbit, say $`p`$-orbit ($`l=1`$). This costs at least another $`\mathrm{}\omega =500`$ MeV for the mass of $`\mathrm{\Theta }^+`$. A question is then whether there is a mechanism to lower the higher state than the negative parity state of $`(0s)^5`$. As shown in Ref. the flavor dependent force due to the Nambu-Goldstone boson exchanges between quarks has a large attraction in the pentaquark state, which compensates the excess of the $`p`$-state energy. Here we illustrate it in the chiral bag model by considering the quark single particle states in a bag as functions of the chiral angle at the bag surface $`F(R)`$ (see Fig. 4) . In the presence of the pion field which interact with the quarks at the bag surface, the equation of motion for the quark field is written as $`(i/{\displaystyle \frac{1}{2}}\mathrm{exp}(i\stackrel{}{\tau }\stackrel{}{\pi }(x)/f_\pi \gamma _5)\delta (rR))\psi =0,`$ (8) where the surface $`\delta `$-function $`\delta (rR)`$ indicates that the interaction occurs at the bag surface. In the hedgehog configuration, $`\stackrel{}{\pi }(x)/f_\pi =\widehat{r}F(r)`$, the quark eigenstates are specified by the parity $`P`$ and the grand spin which is the sum of the orbital angular momentum, spin and isospin, $`\stackrel{}{K}=\stackrel{}{L}+\stackrel{}{S}+\stackrel{}{I}`$. Then, the pion-quark interaction reduces to a spin-isospin interaction of the type $`\stackrel{}{\sigma }\stackrel{}{\tau }`$. For a given $`J=L+S=L\pm 1/2`$, two $`K`$ values are possible, $`K=J\pm 1/2`$. They are degenerate when the pion-quark interaction is zero in the large $`R`$ limit as in the MIT bag model. As the bag radius is reduced and the pion-quark interaction is increased, the degeneracy is resolved. This phenomena is similar to the spin-orbit splitting. In the present case, the state of smaller $`K`$ is lowered, while the other pushed up. The change in the eigenenergies causes level crossing or a rearrangement of the pentaquark state at a certain strength of the pion-quark interaction. A crucial point is that the rearrangement occurs by the crossing of two states of opposite parities, which is followed by a flip of the parity of the pentaquark state. As shown in Fig. 4, for small pion-quark interaction (small $`F(R)`$), the five quark configuration (denoted by the hedgehog quantum numbers $`K^P`$) is $`(0^+)^41^+(0s)^5`$, while it is replaced by the $`(0^+)^41^{}(0s)^41p`$ for $`F(R)\underset{}{>}0.3\pi `$. In the latter, the positive parity state becomes the lowest for the pentaquark state. The parity flip occurs when the pion field is sufficiently strong. As is for the nucleon, if the bag radius of $`\mathrm{\Theta }^+`$ takes a value around $`R0.6`$ where the chiral angle $`F(R)\pi /2`$ , the positive parity $`\mathrm{\Theta }^+`$ can be realized. ### 3.4 Model independent analysis of SU(3) So far, we have discussed models of QCD. Instead, we can work out to a great extent by using only flavor SU(3) symmetry, and derive various relations among masses and coupling constants . Only assumption is that particles of definite spin and parity belong to certain multiplets of SU(3). Symmetry puts constraints on mass and interaction hamiltonians with several parameters, which are determined from experimental data. Since there are more physical quantities than parameters, we can make predictions. If the symmetry is only approximate, we can estimate the breaking effect either by perturbation, or by preparing a wider model space and performing diagonalization. Our interest here is to clarify the nature of $`\mathrm{\Theta }^+`$ and its partners. If SU(3) symmetry is good, they belong to pure $`\overline{10}`$, while if the breaking occurs (from the mass of strange quark), the nucleon and sigma states of the $`\overline{10}`$ start to mix with those of octet members. Presumably, we can imagine that they are also pentaquarks, which however is not a necessary condition . The $`\mathrm{\Xi }`$ states do not mix because of the isospin symmetry. Therefore, we have the following set of particles to consider: $`\mathrm{\Theta }^+(N_8,N_{\overline{10}}),(\mathrm{\Sigma }_8,\mathrm{\Sigma }_{\overline{10}}),\mathrm{\Xi }.`$ We start with writing down the mass matrix for the antidecuplet $`\overline{10}`$ and octet 8, $`H=\left(\begin{array}{cc}M_{\overline{10}}aY& \delta \\ \delta & M_8bY+c[I(I+1)Y^2/4]\end{array}\right),`$ (11) where the parameters are $`M_8,M_{\overline{10}},a,b,c`$ and $`\delta `$, while $`Y`$ and $`I`$ denote hyper charge and isospin. For $`\mathrm{\Theta }^+`$ and $`\mathrm{\Xi }`$ states, only the 11 element is relevant, but for $`N`$ and $`\mathrm{\Sigma }`$ states, the full $`2\times 2`$ matrix must be considered. The six parameters are determined by six inputs from data. Since spin and parity are independent of flavor, we can play with different $`J^P`$’s, and see how the fitting works. We have performed such fittings for $`J^P=1/2^{},1/2^+`$ and $`3/2^{}`$, where there are sufficient number of data. For masses, the three choices of $`J^P`$ work well to a similar extent. The situation, however, changes if the method is applied to decay properties. Since the final state meson and baryon belong to the octet, there are two couplings from $`\overline{10}`$ and $`8`$. Fortunately these two couplings are determined from the decay properties of two known nucleon resonances for the above three cases of $`J^P`$. Using the antidecuplet piece of the coupling constants, we can predict the decay width of $`\mathrm{\Theta }^+`$. The results are summarized in Table 2, where the case of $`1/2^{}`$ is excluded, since it gives too wide widths. Due to ambiguity in the phase of the coupling constants, we have two solutions as listed in the table. From this, we can see that the narrow decay width can be obtained in one solution of $`J^P=3/2^{}`$. This result is natural, since for $`3/2^{}`$, the final $`KN`$ state is d-wave, where the centrifugal barrier suppresses the decay amplitude. The narrow width of $`\mathrm{\Theta }^+`$ could be due to the higher partial wave nature of the decaying channel. ### 3.5 Lattice QCD The investigation of the lattice QCD was started from the early stage of the development . Employing a baryon interpolating field with a suitable five quark configuration, the two point correlation function is studied. The projection into a definite parity state must be also carried out. As inspired by Ref. , this was first performed by Sasaki , who found a resonance-like signal slightly above the $`KN`$ threshold in the $`1/2^{}`$ state. Recently, the significance of the signal has been somewhat weakened, being said that there is no sufficient evidence to deny resonances in the negative parity sector. By now there are several groups who performed simulations, but their results do not always agree with each other and the issue is still controversial . In drawing conclusions, one must know the limitation due to the approximations such as quenched approximation or finite quark mass. Instability of the results depending on the calculation scheme may indicate that the present lattice studies would not be accurate enough for the study of the pentaquark system, or that the pentaquarks might not exist. One of the sources of different results is the use of different types of interpolating field such as: $`J_\mathrm{\Theta }(x)`$ $`=`$ $`ϵ_{abc}[u_a^TC\gamma _5d_b]\{u_e(\overline{s}_e\gamma _5d_c)(ud)\},`$ (12) $`J_\mathrm{\Theta }(x)`$ $`=`$ $`ϵ_{abc}[u_a^TC\gamma _5d_b]\{u_c(\overline{s}_e\gamma _5d_e)(ud)\},`$ (13) $`J_\mathrm{\Theta }(x)`$ $`=`$ $`ϵ_{abc}ϵ_{aef}ϵ_{bgh}[u_e^TCd_f][u_g^TC\gamma _5d_h]C\overline{s}_c^T.`$ (14) Ideally, if computing performance is sufficiently high, the result should not depend on the choice of the interpolating fields. In practice, results depend substantially on the choice. A possibly optimized way is to perform diagonalization of the results of different interpolating fields. Another problem which is physically important is the contamination due to the coupling to the non-resonant scattering state. Recently, an extensive analysis was performed by Takahashi et al. , where they considered a $`2\times 2`$ matrix form of the correlation function generated by the two interpolating fields, (12) and (13), and the matrix was diagonalized to obtain states with an optimal coupling strength. Also, they investigated carefully the volume dependence . They have found a resonance like state which is rather stable against the change of the volume size in the $`1/2^{}`$ sector. They have also studied the spectral weight factors which also supports the resonance like nature of the $`1/2^{}`$ state. However, the resonance signal of the $`1/2^{}`$ channel has been once again questioned in Ref. . Recently, higher spin states have been also investigated . Definitely, further study will be needed to achieve better understanding. ### 3.6 QCD sum rule The QCD sum rule was first applied to the pentaquarks by Zhu , and soon later by Sugiyama et al. with the proper treatment of the parity projection. In this method the two point correlation function for the relevant baryonic state is computed in the asymptotic region in the operator product expansion (OPE). The correlation function in the asymptotic region is then analytically continuated to the low energy region to match the phenomenological spectral function. The method works reasonably well for the ground state baryons and for some resonance states, if the threshold parameter in the phenomenological side is suitably chosen. The validity for excited states is, however, not well tested. The sum rule studies have been performed in many cases for the spin 1/2 sector, and signals of negative parity pentaquarks were seen around 1.5 GeV. Recently, spin 3/2 pentaquarks were also investigated . The observed fact is that the OPE spectral function of the $`1/2^+`$ sector becomes negative or unstable , which indicates either that there is no physical state of $`1/2^+`$ or that the present truncation of the OPE is not good enough. Also, there is a problem of contamination from the $`KN`$ state. In fact, Kondo et al. , claimed the importance of the exclusion of the $`KN`$ scattering state, which may change the result of parity . They have performed the separation by Fierz rearranging the operator of the type of (14) into that of $`KN`$ type. Lee et al. and Kwon et al. also estimated a $`KN`$ component in the two point function by applying the soft kaon theorem . Their estimation showed only a small contribution of the $`KN`$ scattering state to the correlation function and therefore, the result of Sugiyama et. al. is not changed. A fundamental question of the QCD sum rule is its applicability to the pentaquark sector, where the five-quark currents carry higher dimension than the ordinary baryon currents. In such a case, one should include higher orders of OPE in its asymptotic expansion. In this case, however, there emerge more operators of higher dimensions, the vacuum expectation values of which are not known well. ## 4 Decay of $`\mathrm{\Theta }^+`$ Naively, one would expect that the decay of the pentaquark $`\mathrm{\Theta }^+`$ occurs through the fall-apart process as shown in Fig. 5 (left). This should be so in a simple potential model for confinement. However, if, for instance, strings bind quarks, the fall-apart may not necessarily be relevant; if a string breaks by pulling two quarks apart, creation of a quark and antiquark pair must follow. Also, as discussed in Ref. , the fall-apart decay was classified as a forbidden process; it would be interesting if there is such a selection rule which governs some particular decay processes. Here, we adopt the naive picture of the constituent quark model and test the fall-apart process . For mesons and baryons of finite size, antisymmetrization among the four quarks is needed both for the initial state $`\mathrm{\Theta }^+`$ and for the final $`KN`$ scattering states. In the limit of small kaon, however, the exchange term between quarks in the kaon and nucleon can be ignored. In this case, the calculation reduces to the evaluation of the $`\mathrm{\Theta }^+N`$ transition matrix element of the kaon source term, or equivalently of the axial-vector current. In this section we discuss the computation of this process somewhat in detail, since such a method has not been explored much before in hadron physics. ### 4.1 General remark Before going to actual calculations, we briefly look at the general aspect for the width of baryons. Consider a decay of $`\mathrm{\Theta }^+`$ going to the nucleon and kaon. Assuming the spin of the $`\mathrm{\Theta }^+`$, $`J=1/2`$, the interaction lagrangian takes the form $`L_\pm =g_{KN\mathrm{\Theta }}\overline{\psi }_N\gamma _\pm \psi _\mathrm{\Theta }K,`$ (15) where $`\gamma _+=i\gamma _5`$ if the parity of $`\mathrm{\Theta }^+`$ is positive, while $`\gamma _{}=1`$ if the parity of $`\mathrm{\Theta }^+`$ is negative. The decay width is then given by $`\mathrm{\Gamma }_+={\displaystyle \frac{g_{KN\mathrm{\Theta }}^2}{2\pi }}{\displaystyle \frac{M_Nq^3}{E_N(E_N+M_N)M_\mathrm{\Theta }}},\mathrm{\Gamma }_{}={\displaystyle \frac{(E_N+M_N)^2}{q^2}}\mathrm{\Gamma }_+,`$ (16) for the positive ($`+`$) and negative ($``$) parities, where $`M_N`$ and $`M_\mathrm{\Theta }`$ are the masses of the nucleon and $`\mathrm{\Theta }^+`$, and $`E_N=\sqrt{q^2+M_N^2}`$ with $`\stackrel{}{q}`$ being the momentum of the final state kaon in the kaon-nucleon center of mass system, or equivalently in the rest frame of $`\mathrm{\Theta }^+`$. The difference between $`\mathrm{\Gamma }_\pm `$ arises due to the different coupling nature: $`p`$-wave coupling for positive parity $`\mathrm{\Theta }^+`$ and $`s`$-wave coupling for negative parity $`\mathrm{\Theta }^+`$, representing the effect of the centrifugal repulsion in the $`p`$-wave. In the kinematical point of the $`\mathrm{\Theta }^+`$ decay, $`M_\mathrm{\Theta }=1540`$ MeV, $`M_N=940`$ MeV and $`m_K=490`$ MeV, the factor on the right hand side of (16) becomes about 50, which brings a significant difference in the widths of the positive and negative parity $`\mathrm{\Theta }^+`$. If we take $`g_{KN\mathrm{\Theta }}10`$ as a typical strength for strong interaction coupling constants, we obtain $`\mathrm{\Gamma }_+100`$ MeV, while $`\mathrm{\Gamma }_{}`$ 5 GeV. Both numbers are too large as compared with experimentally observed width. Therefore, the relevant question is whether some particular structure of $`\mathrm{\Theta }^+`$ will suppress the above naive values, or not. ### 4.2 Calculation of decay amplitudes The matrix element of a fall-apart decay is written as a product of the spectroscopic factor and an interaction matrix element, $`_{\mathrm{\Theta }^+KN}=S_{KN\mathrm{in}\mathrm{\Theta }^+}h_{int}.`$ (17) The factor $`S_{KN\mathrm{in}\mathrm{\Theta }^+}(\overline{s}q)_K(qqq)_N\|\mathrm{\Theta }^+`$ is a probability amplitude of finding in the pentaquark state three-quark and quark-antiquark clusters having the quantum numbers of the nucleon and kaon, respectively. Calculations of this factor was performed in Refs. . It strongly depends on the internal structure of $`\|\mathrm{\Theta }^+`$. Here we have investigated the four cases; one for the state of $`(0s)^5`$ of $`J^P=1/2^{}`$, and the other three for the $`(0s)^40p`$ of $`J^P=1/2^+`$. The $`(0s)^5`$ configuration is unique, while there are four independent states for the $`(0s)^40p`$ configurations. Among them, we study the one minimizing the spin-flavor interaction of the type $`_{i>j}(\sigma _i\sigma _j)(\lambda _i^f\lambda _j^f)`$ (SF), the one minimizing the spin-color interaction $`_{i>j}(\sigma _i\sigma _j)(\lambda _i^c\lambda _j^c)`$ (SC), and the one of the strong diquark correlations as proposed by Jaffe and Wilczek JW . The resulting spectroscopic factors are summarized in Table 3. Now the interaction matrix element can be computed by the meson-quark interaction of Yukawa type: $`L_{mqq}=ig\overline{q}\gamma _5\lambda _a\varphi ^aq,`$ (18) where $`\lambda _a`$ are SU(3) flavor matrices and $`\varphi ^a`$ are the octet meson fields. The coupling constant $`g`$ may be determined from the pion-nucleon coupling constant $`g_{\pi NN}=5g`$. Therefore, using $`g_{\pi NN}13`$, we find $`g2.6`$. We have calculated the $`\mathrm{\Theta }^+N`$ matrix elements of (18) in the non-relativistic quark model. Further details of calculation can be found in Ref. , and here several results are summarized as follows. For the negative parity state of $`(0s)^5`$, the decay width turns out to be of order of several hundreds MeV or more, typically 0.5 $``$ 1 GeV. In the calculation it has been assumed that the spatial wave function for the initial and final state hadrons are described by a common harmonic oscillator states. Also the masses of the particles are taken as experimental values, e.g., $`M_{\mathrm{\Theta }^+}=1540`$ MeV. For the result of the negative parity state of $`(0s)^5`$, the unique prediction can be made, since there is only one quark model states, meaning that the state can be expressed in terms of the totally antisymmetrized $`KN`$ state. Since there is not a centrifugal barrier, that state is hardly identified with a resonant state with a narrow width. For the positive parity state, we obtain $`\mathrm{\Gamma }=`$ 63 MeV, 32 MeV and 11 MeV, for the SF, SC and JW configurations, respectively. The diquark correlation of (JW) develops a spin-flavor-color wave function having a small overlap with the decaying channel of the nucleon and kaon. In the evaluation of these values, we did not consider spatial correlations. However, if, for instance, small diquarks are developed, spatial overlap becomes less than unity which further suppresses the decay width. In Ref. , such suppression was shown to be significant. However, their absolute values of $`\mathrm{\Gamma }`$ should not be taken seriously, since the PCAC relation was used, which can not be applicable to the quark model calculation. The small values of the decay width for $`J^P=1/2^+`$ as compared with the large values for $`J^P=1/2^{}`$ can be explained by the difference in the coupling structure; one is the pseudoscalar type of $`\stackrel{}{\sigma }\stackrel{}{q}`$ and the other the scalar type of 1. The former of the $`p`$-wave coupling includes a factor $`q/(2M)`$ which suppresses the decay width significantly as compared with the latter at the present kinematics, $`q250`$ MeV and $`M1`$ GeV, when the same coupling constant $`g_{NK\mathrm{\Theta }}`$ is employed. The present analyses can be extended straightforwardly to the $`\mathrm{\Theta }^+`$ of spin 3/2. For the negative parity state, the spin 1 state of the four quarks in the $`\mathrm{\Theta }^+`$ may be combined with the spin of $`\overline{s}`$ for the total spin 3/2. In this case the final $`KN`$ state must be in a $`d`$-wave state, and therefore, the spectroscopic factor of finding a $`d`$-wave $`KN`$ state in the $`(0s)^5`$ is simply zero. If a tensor interaction induces an admixture of a $`d`$-wave configuration, it can decay into a $`d`$-wave $`KN`$ state. However, the mixture of the $`d`$-wave state is expected to be small just as for the deuteron. There could be a possible decay channel of the nucleon and the vector $`K^{}`$ of $`J^P=1^{}`$ . This decay, however, does not occur since the total mass of the decay channel is larger than the mass of $`\mathrm{\Theta }^+`$. Hence the $`J^P=3/2^{}`$ state could be another candidate for the observed narrow state. This state does not have a spin-orbit partner and forms a single resonance peak around its energy. For the positive parity case, the $`p`$-state orbital excitation may be combined with the spin of $`\overline{s}`$ for the total spin 3/2. In this case, the calculation of the decay width is precisely the same as before (See Ref. for more details). After taking the average over the angle $`\stackrel{}{q}`$, however, the coupling yields the same factor as for the case $`J=1/2`$. Hence the decay rate of spin 3/2 $`\mathrm{\Theta }^+`$ is the same as that of $`\mathrm{\Theta }^+`$ of spin 1/2 in the present treatment, if the mass of the $`3/2^+`$ state is the same as the $`1/2^+`$ state. ## 5 Production of $`\mathrm{\Theta }^+`$ The $`\mathrm{\Theta }^+`$ production from the non-strange initial hadrons is furnished by the creation of $`s\overline{s}`$ pair, which requires energy deposit of around 1 GeV. In general, the reaction mechanism of such energy region is not well understood. However, as one of practical methods, we adopt an effective lagrangian approach and perform computations of Born (tree) diagrams. Input parameters in the lagrangians reflect the properties of $`\mathrm{\Theta }^+`$ and therefore, the comparison of calculations and experiments will help study the structure of $`\mathrm{\Theta }^+`$. In particular, extraction of spin and parity is the important purpose in the study of reactions. Here we briefly discuss (1) photoproduction as originally performed in experiment by the LEPS group , and (2) $`\mathrm{\Theta }^+\mathrm{\Sigma }^+`$ production induced by the polarized $`\stackrel{}{p}\stackrel{}{p}`$ for the determination of the parity . ### 5.1 Photoproduction As described in detail in Ref. , in the effective lagrangian method we calculate the Born (tree) diagrams as depicted in Fig. 6. The actual form of the interaction lagrangian depends on the interaction schemes, i.e., either pseudoscalar (PS) or pseudovector (PV). In the PS, the three Born diagrams (a)-(c) are computed with the gauge symmetry maintained. In the PV, on the contrary, the contact Kroll-Ruderman term (d) is also necessary. In the PS scheme, the contact term may be included in the antinucleon contribution of the nucleon Z-diagram. If chiral symmetry is respected, the low energy theorem guarantees that the two schemes provide the same result in the low energy limit. In reality, due to the large energy deposit of order 1 GeV, the equivalence is violated. It is shown that the difference in the two schemes is proportional to the photon momentum in the first power (which therefore vanishes in the low energy limit) and to the anomalous magnetic moment of $`\mathrm{\Theta }^+`$ . In order to express the finite size effect of the nucleon, we need to consider the form factor. Here we adopt a gauge invariant one with a four momentum cutoff . This form factor suppresses the nucleon pole contributions in the PS scheme (and hence the contact term also in the PV scheme), as reflecting the fact that the nucleon intermediate state is far off-shell. Consequently, the dominant contribution is given by the t-channel process of the kaon exchange and/or $`K^{}`$ meson exchange. The ambiguity of the anomalous magnetic moment of $`\mathrm{\Theta }^+`$ is also not important. Therefore, the difference between the PS and PV schemes is significantly suppressed when kaon exchange term is present as for the case of the neutron target. This allows one to make rather unambiguous theoretical predictions. We have computed the photoproduction of $`\mathrm{\Theta }^+`$ from the neutron and proton, and first for $`J^P=1/2^\pm `$. Here are several remarks: 1. When the decay width $`\mathrm{\Gamma }_{\mathrm{\Theta }^+KN}=15`$ MeV is used the typical total cross section values are about 100 \[nb\] for the positive parity and about 10 \[nb\] for the negative parity. Since the total cross section is proportional to the decay width $`\mathrm{\Gamma }_{\mathrm{\Theta }^+KN}`$, experimental information on the decay width is important to determine the size of cross sections, or vise versa. For instance, for the decay width about 1 MeV or less, the total cross sections will be of order of 10 \[nb\] or less for the positive parity and 1 \[nb\] or less for the negative parity. In general the cross sections are about ten times larger for the positive parity $`\mathrm{\Theta }^+`$ than for the negative parity. The $`p`$-wave coupling $`\stackrel{}{\sigma }\stackrel{}{q}`$ effectively enhances the coupling strength by factor 3 – 4 as compared with the $`s`$-wave coupling for the negative parity, when the momentum transfer amounts to 1 GeV. 2. For the neutron target, the kaon exchange term is dominant. In this case, the $`K^{}`$ contributions are not important even with a large $`K^{}N\mathrm{\Theta }`$ coupling $`\|g_{K^{}N\mathrm{\Theta }}\|=\sqrt{3}\|g_{KN\mathrm{\Theta }}\|`$ . Hence the theoretical prediction for the neutron target is relatively stable. The angular dependence has a peak at $`\theta 60`$ degrees in the center-of-mass system, a consequence of the vertex structure of the $`\gamma KK`$ vertex in the kaon exchange term. Since this feature is common to both parities, the difference in the parity of $`\mathrm{\Theta }^+`$ may not be observed in the angular distribution. 3. The kaon exchange term vanishes for the case of the proton target. Therefore, the amplitude is a coherent sum of various Born terms, where the role of the $`K^{}`$ exchange is also important. The theoretical prediction for the proton target is therefore rather difficult. Very recently, the CLAS collaboration has reported no significant evidence of $`\mathrm{\Theta }^+`$ in the reaction $`\gamma p\overline{K}^0\mathrm{\Theta }^+`$ . This result should be taken seriously, because they have achieved significantly higher statistics as compared to the previous experiments performed close to the threshold. What could then be the fate of $`\mathrm{\Theta }^+`$? Yet their result does not lead to the absence of $`\mathrm{\Theta }^+`$ immediately, because the previous positive evidences were seen mostly in the reactions from the neutron. Due to the violation of isospin symmetry in the electromagnetic interaction, there could be asymmetry in the reactions from the proton and neutron. A well-known example is the Kroll-Ruderman term in the pion photoproduction, which survives only in the charge exchange channels such as $`\gamma p\pi ^+n`$. We have performed a calculation using once again an effective Lagrangian, but with $`J^P=3/2^\pm `$ $`\mathrm{\Theta }^+`$ . There is a significant difference between $`1/2`$ and $`3/2`$ cases; only PV formalism is possible for the latter case. Without the equivalence between the PS and PV, the role of the contact term in the PV scheme is very much different for the two spin cases. In the $`3/2`$ case, the contact term dominates and the production rate from the neutron is large but that from the proton is strongly suppressed. Using the decay width $`\mathrm{\Gamma }1`$ MeV, the cross section of $`\gamma p\overline{K}^0\mathrm{\Theta }^+`$ was estimated to be a few nb which does not contradict the CLAS data. Further information from the neutron target is crucially important to settle the problem of the pentaquarks. ### 5.2 Polarized proton beam and target This reaction was considered in order to determine the parity of $`\mathrm{\Theta }^+`$ unambiguously, independent of any reaction mechanisms . In the photoproduction case, the determination of parity is also possible if we are able to control the polarization of both the initial and the final states , which is however very difficult in the present experimental setup. The system of two protons provides a selection rule due to Fermi statistics. Since the isospin is $`I=1`$, the spin and angular momentum of the initial state must be either $`(S,L)=(0,\mathrm{even})`$ or $`(S,L)=(1,\mathrm{odd})`$. Now consider the reaction $`\stackrel{}{p}+\stackrel{}{p}\mathrm{\Theta }^++\mathrm{\Sigma }^+.`$ (19) at the threshold region, where the relative motion in the final state is in $`s`$-wave. It is shown that if the initial spin state has $`S=0`$, then the parity of the final state is positive and hence the parity of $`\mathrm{\Theta }^+`$ MUST BE positive. Likewise, if $`S=1`$ the parity of $`\mathrm{\Theta }^+`$ MUST BE negative. This idea is similar to the one used to determine the parity of the pion . One can compute production cross sections by employing an effective lagrangian of the kaon and $`K^{}`$ exchange model. The strength of the $`KN\mathrm{\Theta }^+`$ vertex is determined once again for $`\mathrm{\Gamma }=15`$ MeV. The $`K^{}N\mathrm{\Theta }^+`$ vertex is expressed as the sum of the vector and tensor terms. The vector and tensor couplings are unknown, but here we take their strengths to be $`\|g_{K^{}N\mathrm{\Theta }}^V\|=(1/2)\|g_{K^{}N\mathrm{\Theta }}^T\|=\|g_{KN\mathrm{\Theta }}\|`$. The signs are then tested for all four possible cases to see the effect of the couplings. For the $`KN\mathrm{\Sigma }`$ and $`K^{}N\mathrm{\Sigma }`$ coupling, we take the phenomenological one from the Nijmegen potential . The monopole form factor is then introduced at each vertex, with the same cutoff parameter $`\mathrm{\Lambda }=1`$ GeV for simplicity. The choice of the interaction parameter is important for such high momentum transfer reactions. It should reflect the structure of the nucleon which has the size of order of 0.5 fm or larger. Hence $`\mathrm{\Lambda }1`$ GeV is crudely the upper limit which is compatible with nucleon size $`\underset{}{<}0.5`$ fm. In fact, the parameters of the Nijmegen soft core potential are chosen by such a consideration. The results are shown in Fig. 7 for both positive and negative parity $`\mathrm{\Theta }^+`$. The selection rule for the positive and negative parity $`\mathrm{\Theta }^+`$ is shown clearly by the energy dependence at the threshold region. For the allowed channel the final state is in an $`s`$-wave with the energy dependence from the threshold $`(EE_{th})^{1/2}`$, whereas for the forbidden channel the partial wave of the final state is $`p`$-wave with the energy dependence $`(ss_{th})^{3/2}`$. Recently COSY-TOF reported the result for the $`\mathrm{\Theta }^+\mathrm{\Sigma }^+`$ production in the unpolarized $`pp`$ scattering at $`p_p=2.95`$ GeV/c . They quote the total cross section $`\sigma 0.4\pm 0.2`$ \[$`\mu `$b\] at 30 Mev above the threshold in the center of mass energy. In comparison with theory, if we adopt a narrower width of about 5 MeV, the cross section will be about 0.5 \[$`\mu `$b\] for the positive parity and 0.05 \[$`\mu `$b\] for the negative parity. This comparison seems to favor the positive parity $`\mathrm{\Theta }^+`$. Although there remain some ambiguities in theoretical calculations, such a comparison of the total cross section will be useful to distinguish the parity of $`\mathrm{\Theta }^+`$. An alternative quantity which is powerful for the determination of the parity is the spin polarized quantity as defined by $`A_{xx}=(^3\sigma _0+^3\sigma _1)/(2\sigma _0)1,`$ (20) where $`\sigma _0`$ and $`\sigma _1`$ are the cross sections for the spin singlet and triplet states of the two protons. By taking the ratio of the two cross sections, ambiguities of various coupling constants and form factors are nearly cancelled. The observation of $`A_{xx}`$ as well as the energy dependence of the polarized cross sections will provide the best opportunity to determine the parity of $`\mathrm{\Theta }^+`$ . ## 6 Summary In this note we have discussed several aspects of the pentaquark baryon $`\mathrm{\Theta }^+`$ including its structure and production reactions. Among various properties of $`\mathrm{\Theta }^+`$, the importance of the spin and parity in relation with the decay width has been emphasized. The relevant points are as follows: 1. The chiral force may change the quark energy levels; with a sufficient strength, an $`l=1`$ orbit may be lower than the $`l=0`$ orbit. Consequently, a positive parity state can be the lowest pentaquark state. Better understanding of the role of the Nambu-Goldstone bosons is very important. 2. The decay of the pentaquark state through the fall-apart process is sensitive to the internal quark structure, especially to the spin and parity. It was shown that for $`J^P=1/2^{}`$, the naive ground state of $`(0s)^5`$ can no longer survive as a narrow resonance as the decay width is unphysically large. In fact, the uniqueness of the $`(0s)^5`$ configuration implies that that configuration can be written in terms of the $`KN`$ state which can no longer be bound in the naive quark model where the confining force vanishes between the two color singlet states. Contrary, the decay widths of $`1/2^+`$ states were obtained to be of order ten MeV, but with once again strong dependence on the configuration. We have also commented on the possibility of higher spin states, especially $`J=3/2`$. The $`J^P=3/2^{}`$ state could be an interesting alternative possibility. 3. As recently pointed by Hiyama , the coupling to the $`KN`$ decay channel is extremely important when considering the five-body system seriously; it could change the nature of a confined pentaquark configuration completely. Since the five quark configuration must have components of two color singlet hadrons, one (or some) of them can be a decay channel(s), unless there are particular selection rules. The $`J^P=1/2^\pm `$ states are not subject to any such selection rules and must be accompanied by such a coupling to the two hadron ($`KN`$) scattering state. In such a case, a coupled channel treatment is mandatory. In contrast, the $`J^P=3/2^{}`$ state, unless there is $`d`$-wave mixing with the $`(0s)^5`$ configuration, the angular momentum conservation forbids the coupling of the state with the decay channel. 4. So far, our understanding of the reaction mechanism for the pentaquark production is rather limited. Perhaps, the best we can do is to use an effective lagrangian approach with a reasonable choice of model parameters. 5. In the photoproduction of $`\mathrm{\Theta }^+`$, it was found a large asymmetry between the reactions from the proton and neutron targets, especially when $`J=3/2`$ . In relation with the understanding of the narrow decay width, we once again mention that the higher spin state would be an interesting possibility. 6. In order to determine the parity of the pentaquark, the polarized proton scattering $`\stackrel{}{p}\stackrel{}{p}\mathrm{\Theta }^+\mathrm{\Sigma }`$ provides a model independent method. Measurement of such reaction is extremely important to further explore the physics of pentaquarks. The current situation for the pentaquarks is not settled. Of course, the experimental confirmation is the most important issue. However, we have also seen that different theoretical approaches make different predictions. These facts imply that there could be more aspects that we do not know yet well about the low energy QCD. Perhaps the pentaquark has provided us with an ideal opportunity to explore further challenges to the problems. ### Acknowledgements The author would like to thank the hospitality to the organizers of the workshop on HADRON PHYSICS, March 7 – 17, (2005) Puri, India, during his stay. He thanks K. Hicks, E. Hiyama, T. Hyodo, M. Kamimura, H.C. Kim, T. Nakano, S.I. Nam, M. Oka, E. Oset, A. Titov, H. Toki, A.W. Thomas and M.J. Vicente-Vacas for discussions and collaborations. This work supported in part by the Grant for Scientific Research ((C) No.16540252) from the Ministry of Education, Culture, Science and Technology, Japan.
warning/0506/quant-ph0506189.html	ar5iv	text	# General paradigm for distilling classical key from quantum states ## I Introduction We often want to communicate with friends or strangers in private. Classically, this is impossible if we wish to communicate over long distances, unless we have met before with our friend and exchanged a secret key which is as long as the message we want to send. On the other hand, quantum cryptography allows two people to communicate privately with only a very short key which is just used to authenticate the message. Every quantum cryptographic protocol is equivalent to the situation where both parties (Alice and Bob) share some quantum state $`\rho _{AB}`$, and then perform local operations on that state and engage in public communication (LOPC) to obtain a key which is private from any eavesdropper. Until recently, every quantum protocol was also equivalent to distilling pure entanglement from this shared state. I.e., achieving privacy was equivalent to the two parties converting many copies of the state $`\rho _{AB}`$, to a smaller number of pure EPR pairs $$\|\psi _+=\frac{1}{\sqrt{2}}(\|00+\|11),$$ (1) using local operations and classical communication (LOCC), and then performing a measurement on the EPR pairs in the computational basis. Examples of such protocols include BB84 , B92 , and of course, E91 . It was thus thought that achieving security is equivalent to distilling pure entanglement, and a number of results pointed in this direction . Recently, however, we have shown that this is not the case – there exist examples of bound entangled states which can be used to obtain a secret key . Bound entangled states are ones which need pure entanglement to create, but no pure entanglement can be distilled from them. This helps explain the properties of bound entangled states. They have entanglement which protects correlations from the environment (or an eavesdropper), but the entanglement is so twisted that it can’t be brought into pure form. This then raised the question of what types of quantum states provide privacy. In we were able to find the general form of private quantum states $`\gamma _{ABA^{}B^{}}`$. This allowed us to recast the theory of privacy (under local operations and public communication – or LOPC) in terms of entanglement theory (local operations and classical communication – or LOCC). In entanglement theory, the basic unit is the EPR pair, while in privacy theory, the only difference is that one replaces the EPR pair with general private states $`\gamma _{ABA^{}B^{}}`$ as the basic units. In the present article, we review the results of in greater detail, and expand on the proofs and tools. Namely, we study and show that the general form of a private state on a Hilbert space $`_A_A^{}_B_B^{}`$ with dimensions $`d_A=d_Bd`$, $`d_A^{}`$ and $`d_B^{}`$, is of the form $$\gamma _{ABA^{}B^{}}=UP_{AB}^+\sigma _{A^{}B^{}}U^{},$$ where $`P_{AB}^+`$ is a projector onto the maximally entangled state $`\psi _+=_i\frac{1}{\sqrt{d}}\|e_if_i`$, and $`U`$ is the arbitrary twisting operation $$U=\underset{k,l=0}{\overset{d1}{}}\|e_kf_le_kf_l\|_{AB}U_{A^{}B^{}}^{kl}.$$ The key is obtained after measuring in the $`\|e_if_i`$ basis. We will henceforth refer to $`\psi _+`$ as the maximally entangled or EPR state (or Bell state in dimension $`2\times 2`$). We show that the rate of key $`K_D`$ which can be obtained from a quantum state can be strictly greater than the distillable entanglement, and this even holds if the distillable entanglement is strictly zero. We also show that the size of the private key is generally bounded from above by the regularized relative entropy of entanglement $`E_r^{\mathrm{}}`$ . This will be sufficient to prove that one can have a maximal rate of key strictly less than the entanglement cost (the number of singlets required to prepare a state under LOCC). In section II we introduce some of the basic concepts and terminology we will use throughout the paper. This includes the notion of private states, pbits which contain one bit of private key, and pdits which have many bits of key. In Section III we show that a state is secure if and only if it is of the form given above. Then we show different useful ways to write the private states in Section IV, and give some useful examples, and examine some of their properties. This includes the notion of irreducibility which is used to define the basic unit of privacy for private states. States which have a perfect bit of key must have some distillable entanglement . The case of bound entangled states with secure key is only found in the case of states which are not perfectly secure, although they are arbitrarily secure. This motives our investigation in Section V of approximate pbits. We then demonstrate how to rewrite a bipartite state in terms of the eavesdropper’s density matrix in Section VI. This allows us to interpret previous results in terms of the eavesdropper’s states. Then, in Section VII, we summarize the previous results in preparation for showing that bound entangled states can have a key. In Sections VIII and IX, we review the paradigms of entanglement (LOCC) and privacy theory (LOPC), and show the equivalence of key rates in the two paradigms. We then discuss and compare security criteria in these paradigms in Section XV-C of the Appendix. In Section X we give a number of bound entangled states and show that they can produce a private key. The methods allow one to find a wide class of states which are bound entangled, because the fact that they have key automatically ensures that they are entangled, which is usually the difficult part in showing that a state is bound entangled (the PPT criteria can be quickly checked to see that the states are non-distillable). In Section XI we prove that the relative entropy distance is an upper bound on the rate of key. In Section XII, a class of NPT states are introduced which appear to be bound entangled. They are derived from a class of bound entangled private key states. An additional result discussed in Section XIII, which we only mentioned in passing in , is that the bound entangled key states can be used as the basis of a cryptographic primitive we call a controlled private quantum channel. We conclude in Section XIV with a few open questions. ## II Security contained in quantum states In this section we will introduce the class of states which contain at least one bit (or dit) of perfectly secure key which is directly accessible – these we call private bits (or dits). We discuss the properties of these states and argue the generality of this approach. In particular we introduce the notion of twisting, which is a basic concept in dealing with private states. A well known state that contains one bit of secure key which is directly accessible is the singlet state. After measuring it in a local basis Alice and Bob obtain bits that are perfectly correlated with each other and completely uncorrelated with the rest of the world including an eavesdropper Eve. This is because the singlet state as a whole is decoupled from the environment, being a pure state. However even if a state is mixed, it can contain secure key. Yet the key must then be located only in a part of it. More formally, we consider a four-partite mixed state $`\rho _{ABA^{}B^{}}`$ of two systems $`A,A^{}`$ belonging to Alice and $`B,B^{}`$ belonging to Bob. The $`AB`$ subsystem of the state will be called the key part of the state – it is the part of the state which produces key upon measurement. The $`A^{}B^{}`$ subsystem will be called the shield of the state. It is called this, because its presence is what will cause the $`AB`$ part of the state to be secure, by shielding information from an eavesdropper. We assume the worst case scenario – that the state is the reduced density matrix of the pure state $`\psi _{ABA^{}B^{}E}`$ where we trace out the system $`E`$ belonging to eavesdropper Eve. We then distinguish a product basis $`=\{e_i,f_j\}`$ in system $`AB`$. For our purposes, without loss of generality, we often choose $``$ to be the standard basis $`\{\|ij\}`$. Distinguishing the basis is connected with the fact that we are dealing with classical security, which finally is realized in some fixed basis. Now, consider the state of systems $`ABE`$ after measurement performed in the basis $``$ by Alice and Bob. This state is of the form $$\rho _{ccq}=\underset{i,j=0}{\overset{d1}{}}p_{ij}\|e_if_j_{AB}e_if_j\|\rho _{ij}^E.$$ (2) The above form of the state is usually called a ccq state. We will therefore refer to a ccq state associated with state $`\rho _{ABA^{}B^{}}`$, and it is understood that it is also related to chosen basis $``$. The distribution $`p_{ij}`$ will sometimes be referred to as the distribution of the ccq state. We can now distinguish types of states $`\rho _{ABA^{}B^{}}`$ via looking at their ccq states (always assuming that some fixed basis $``$ was chosen): ###### Definition 1 A state $`\rho _{ABA^{}B^{}}`$ is called secure with respect to a basis $`\{\|e_if_j_{AB}\}_{i,j=1}^d`$ if the state obtained via measurement on AB subsystem of its purification in basis $``$ followed by tracing out $`A^{}B^{}`$ subsystem (i.e. its ccq state) is product with Eve’s subsystem: $$\left(\underset{i,j=0}{\overset{d1}{}}p_{ij}\|e_if_je_if_j\|_{AB}\right)\rho _E.$$ (3) Such a state $`\rho _{ABA^{}B^{}}`$ will be also called ”$``$ secure”. Moreover if the distribution $`\{p_{ij}\}=\{\frac{1}{d}\delta _{ij}\}`$ so that the ccq state is of the form $$\left(\underset{i=0}{\overset{d1}{}}\frac{1}{d}\|e_if_ie_if_i\|_{AB}\right)\rho _E,$$ (4) the state $`\rho _{ABA^{}B^{}}`$ is said to have $``$-key. One can ask when two states $`\rho _{ABA^{}B^{}}`$ and $`\sigma _{ABA^{}B^{}}`$ are equally secure with respect to a given product basis $``$. First let us define what does it mean ”equally secure”. A natural definition would be that when Alice and Bob measure systems AB in the basis, then Eve by any means cannot distinguish between two situations, as far as the outcomes of the measurement are concerned. In particular, the states are definitely equally secure, when their ccq states are equal. For our purpose we will need to know when for two states the latter relation holds. It is obvious that any unitary transformation applied to systems $`A^{}B^{}`$ of the state $`\rho _{ABA^{}B^{}}`$ will not change the ccq state. (Note that it cannot be just any CP map; for example, partial trace of systems $`A^{}B^{}`$ would mean giving it to Eve, which of course would change the ccq states). As it will be demonstrated in the next section, we can actually do much more without changing the ccq state. Namely we can apply an operation called ”twisting”. This operation is defined for system $`ABA^{}B^{}`$ and with respect to a product basis $``$ of $`AB`$ system as follows: ###### Definition 2 Given product basis $`=\{e_i,f_j\}_{k,l}`$ on systems $`AB`$, the unitary operation acting on system $`ABA^{}B^{}`$ of the form $$U=\underset{k,l=0}{\overset{d1}{}}\|e_kf_le_kf_l\|_{AB}U_{A^{}B^{}}^{kl},$$ (5) is called $``$-twisting, or shortly twisting. Finally we define the class of private states. The states from that class are proven to be the only quantum states which after measurement on Alice and Bob subsystems give an ideal key. In other words these are the only states from which Alice and Bob can get an ideal ccq state (4) according to definition 1 of security. For the sake of clarity, we recall this proof with details in Section III. ###### Definition 3 A state $`\rho _{ABA^{}B^{}}`$ of a Hilbert space $`_A_A^{}_B_B^{}`$ with dimensions $`d_A=d_Bd`$, $`d_A^{}`$ and $`d_B^{}`$, of the form $$\gamma ^{(d)}=\frac{1}{d}\underset{i,j=0}{\overset{d1}{}}\|e_if_ie_jf_j\|_{AB}U_i\sigma _{A^{}B^{}}U_j^{},$$ (6) where the state $`\sigma _{A^{}B^{}}`$ is an arbitrary state of subsystem $`A^{}B^{}`$, $`U_i`$’s are arbitrary unitary transformations and $`\{\|e_i\}_{i=0}^{d1}`$, $`\{\|f_j\}_{i=0}^{d1}`$ are local basis on $`_A`$ and $`_B`$ respectively, is called private state or pdit. In case of $`d=2`$ the state is called pbit. Note, that maximally entangled states are also private states, which is in case when $`d_A^{}=d_B^{}=1`$. In general, any pdit can be created out of a maximally entangled state with additional state on $`\sigma _{A^{}B^{}}`$ (which we will call basic pdit) by some twisting. ###### Definition 4 A state $`\rho _{ABA^{}B^{}}`$ of a Hilbert space $`_A_A^{}_B_B^{}`$ with dimensions $`d_A=d_Bd`$, $`d_A^{}`$ and $`d_B^{}`$, of the form $$\rho _{ABA^{}B^{}}=P_{AB}^+\sigma _{A^{}B^{}},$$ (7) is called a basic pdit. ###### Remark 1 Let us note, that one could define states with ideal key also in a different way than in definition 1. Namely, we could say that the state $`\rho _{AB}`$ has $``$ key iff it has a subsystem $`ab`$ such, that the subsystem $`abE`$ of its purification $`\|\psi _{ABE}`$ is an ideal ccq state of eq. (4), where $`=\{\|e_if_i\}_{i=1}^d`$. However maximally entangled state is not of this form, but is locally equivlent to a state of this form, i.e. they would be transformable into one another by means of unitary embeddings or partial isometries. In fact the whole class of so defined states with ideally secure key, would be locally equivalent to the one we have introduced in definition 3, and by characterization of the latter, equivalent to the class of private states. Another definition of states with ideal key could be as follows. A state $`\rho _{AB}`$ is called to have key if there are operations $`\mathrm{\Lambda }_A`$ and $`\mathrm{\Lambda }_B`$ with Kraus operators $`\{\|i_{K_A}X_A^{(i)}\}`$ and $`\{\|i_{K_B}Y_B^{(i)}\}`$ respectively, such that the $`ccq`$ state of subsystems $`K_AK_BE`$ of the purification of an output state $`\mathrm{\Lambda }_A\mathrm{\Lambda }_B(\rho _{AB})`$, is an ideal ccq state of eq. (4). This definition however would not allow for easy characterization of this class of states, and still, such states would be locally equivalent to private states defined in 3. ### II-A Some facts and notations In what follows, by $`\|\|.\|\|`$ we mean the trace norm, i.e. the sum of the singular values of an operator. For any bipartite operator $`XB(_1_2)`$, by $`X^\mathrm{\Gamma }`$ we mean the partial transposition of $`X`$ with respect to system 2, that is: $$(I_1T_2)X,$$ (8) where $`T_2`$ denotes the matrix transposition over system $`2`$ of a matrix $`X`$. For brevity, we will use the same symbol for any partial transposition. In particular, we deal often with systems of four subsystems $`ABA^{}B^{}`$, and we will take partial transposition with respect to subsystems $`B`$ and $`B^{}`$. Hence $`\rho _{ABA^{}B^{}}^\mathrm{\Gamma }`$ denotes $`(I_AT_BI_A^{}T_B^{})(\rho _{ABA^{}B^{}})`$, with $`T_B`$ and $`T_B^{}`$ denoting the matrix transposition on systems $`B`$ and $`B^{}`$ respectively. To give example of a mixed notation, we consider $`\rho _{ABA^{}B^{}}B(𝒞^2𝒞^2𝒞^d𝒞^d)`$. In block matrix form, such a state has a bipartite structure of blocks, so that it reads: $$\rho _{ABA^{}B^{}}=\left[\begin{array}{cccc}A_{0000}& A_{0001}& A_{0010}& A_{0011}\\ A_{0100}& A_{0101}& A_{0110}& A_{0111}\\ A_{1000}& A_{1001}& A_{1010}& A_{1011}\\ A_{1100}& A_{1101}& A_{1110}& A_{1111}\end{array}\right].$$ (9) This state after partial transposition with respect to system $`BB^{}`$ reads: $$\rho _{ABA^{}B^{}}^\mathrm{\Gamma }=\left[\begin{array}{cccc}A_{0000}^\mathrm{\Gamma }& A_{0100}^\mathrm{\Gamma }& A_{0010}^\mathrm{\Gamma }& A_{0110}^\mathrm{\Gamma }\\ A_{0001}^\mathrm{\Gamma }& A_{0101}^\mathrm{\Gamma }& A_{0011}^\mathrm{\Gamma }& A_{0111}^\mathrm{\Gamma }\\ A_{1000}^\mathrm{\Gamma }& A_{1100}^\mathrm{\Gamma }& A_{1010}^\mathrm{\Gamma }& A_{1110}^\mathrm{\Gamma }\\ A_{1001}^\mathrm{\Gamma }& A_{1101}^\mathrm{\Gamma }& A_{1011}^\mathrm{\Gamma }& A_{1111}^\mathrm{\Gamma }\end{array}\right].$$ (10) In the above equation, the partial transposition on LHS is with respect to system $`BB^{}`$ and on the RHS, only with respect to system $`B^{}`$, as the partial transposition with respect to system $`B`$, which is a one qubit system, resulted already in appropriate reordering of the block operators $`A_{ijkl}`$. In what follows, we will repeatedly use equivalence of the trace norm distance and fidelity proved by Fuchs and van de Graaf : ###### Lemma 1 For any states $`\rho `$, $`\rho ^{}`$ there holds $$1F(\rho ,\rho ^{})\frac{1}{2}\rho \rho ^{}\sqrt{1F(\rho ,\rho ^{})^2}$$ (11) Here $`F(\rho ,\rho ^{})=\mathrm{Tr}\sqrt{\sqrt{\rho }\rho ^{}\sqrt{\rho }}`$ is fidelity; We also use the Fannes inequality (in the form of ): $$\|S(\rho )S(\sigma )\|2\rho \sigma \mathrm{log}d+h(\rho \sigma ),$$ (12) which holds for arbitrary states $`\rho `$ and $`\sigma `$ satisfying $`\rho \sigma 1`$. ### II-B On twisting and privacy squeezing Here we will show that twisting does not change the ccq state arising from measurement of the key part. Then we will introduce a useful tool by showing that twisting can pump entanglement responsible for security of ccq state into the key part. We have the following theorem. ###### Theorem 1 For any state $`\rho _{AA^{}BB^{}}`$ and any $``$-twisting operation $`U`$, the states $`\rho _{AA^{}BB^{}}`$ and $`\sigma _{ABA^{}B^{}}=U\rho _{AA^{}BB^{}}U^{}`$ have the same ccq states w.r.t $``$, i.e. after measurement in basis $``$, the corresponding ccq states are equal: $`\stackrel{~}{\rho }_{ABE}=\stackrel{~}{\sigma }_{ABE}`$ ###### Proof: To show that subsystem $`\rho _{ABE}`$ is not affected by $``$ controlled unitary with a target on $`A^{}B^{}`$ we will consider the whole pure state: $$\|\psi _{ABA^{}B^{}E}=\underset{ijklm}{}a_{ijklm}\|ijklm\|\psi $$ (13) (without loss of generality we take $``$ to be standard basis). After von Neumann measurement on $``$ and tracing out the $`A^{}B^{}`$ part the output state is the following: $$\stackrel{~}{\rho }_{ABE}=\underset{ijklmn}{}a_{ijklm}\overline{a}_{ijkln}\|ijij\|\|mn\|.$$ (14) Let us now subject $`\|\psi `$ to controlled unitary $`U_{ABA^{}B^{}}I_E`$, $$U_{ABA^{}B^{}}I_E\|\psi =\underset{ijklm}{}a_{ijklm}\|ijU^{ij}\|kl\|m\|\stackrel{~}{\psi },$$ (15) and then on the output state $`\|\stackrel{~}{\psi }`$ perform a complete measurement on $``$ reading the output: $`P_{ij}\|\stackrel{~}{\psi }\stackrel{~}{\psi }\|P_{ij}={\displaystyle \underset{klmstn}{}}a_{ijklm}\overline{a}_{ijstn}`$ $`\|ijij\|_{AB}U^{ij}\|klst\|(U^{ij})_{A^{}B^{}}^{}\|mn\|_E.`$ (16) Performing partial trace and summing over $`i,j`$ we obtain the same density matrix as in (14) which ends the proof. The above theorem states that two states which differ by some twisting $`U`$, have the same $`ccq`$ state obtained by measuring their key parts, and tracing out their shields. However, since twisting does not affect only the ccq state, one can be interested in how the whole state changes when subjected to such an operation. We will show now an example of twisting which will be of great importance for further considerations in this paper. Subsequently, we will construct from this twisting an operation called privacy squeezing (shortly: p-squeezing), which shows the importance of the above theorem. The operation of p-squeezing is a kind of primitive in the paradigm which we will present in the paper. Consider the following technical lemma: ###### Lemma 2 For any state $`\sigma _{ABA^{}B^{}}(𝒞^2𝒞^2𝒞^d𝒞^d^{})`$ expressed in the form $`\sigma _{ABA^{}B^{}}=_{ijkl=0}^1\|ijkl\|A_{ijkl}`$ there exists twisting $`U_{ps}`$ such that if we apply this to $`\sigma _{ABA^{}B^{}}`$, and trace out $`A^{}B^{}`$ part, the resulting state on $`AB`$ $`\rho _{AB}=\mathrm{Tr}_{A^{}B^{}}[U_{ps}\sigma _{ABA^{}B^{}}U_{ps}^{}]`$ will have the form $$\rho _{AB}=\left[\begin{array}{cccc}\times & \times & \times & A_{0011}\\ \times & \times & \times & \times \\ \times & \times & \times & \times \\ \times & \times & \times & \times \end{array}\right],$$ (17) where we omit non-important elements of $`\rho _{AB}`$. ###### Proof: Twisting, by its definition (5) is determined by the set of unitary transformations. In the case of pbit which we now consider, there are four unitary transformations which determine it: $`\{U_{kl}\}_{k,l=0}^1`$. Let us consider singular value decomposition of the operator $`A_{0011}`$ to be $`VR\stackrel{~}{V}`$ with $`V,\stackrel{~}{V}`$ unitary transformations, and $`R`$ \- nonnegative diagonal operator. Note, that by unitary invariance of norm, we have that $`A_{0011}=R=\mathrm{Tr}R`$. We then define twisting by choosing $`U_{00}=V^{}`$, $`U_{11}=\stackrel{~}{V}`$, and $`U_{01}=U_{10}=I`$. The $`AB`$ subsystem of twisted $`\sigma _{ABA^{}B^{}}`$ state is $$\rho _{AB}=\underset{ijkl=0}{\overset{1}{}}\mathrm{Tr}(U_{ij}A_{ijkl}U_{kl}^{})\|ijkl\|,$$ (18) so for such chosen twisting we have indeed, that the element $`\|0011\|`$ of the matrix of $`\rho _{AB}`$ is equal to $`\mathrm{Tr}U_{00}^{}VR\stackrel{~}{V}U_{11}^{}=\mathrm{Tr}R=A_{0011}`$, which proves the assertion. We will give now the following corollary, which will serve as simple exemplification of this result. ###### Corrolary 1 Let the key part be two qubit system. Consider then a state of the form (where blocks are operator acting on $`A^{}B^{}`$ system): $$\sigma _{ABA^{}B^{}}=\left[\begin{array}{cccc}A_{0000}& 0& 0& A_{0011}\\ 0& A_{0101}& A_{0110}& 0\\ 0& A_{1001}& A_{1010}& 0\\ A_{1100}& 0& 0& A_{1111}\end{array}\right],$$ (19) there exists twisting such that the state after partial trace on $`A^{}B^{}`$ has a form $$\rho _{AB}=\left[\begin{array}{cccc}A_{0000}& 0& 0& A_{0011}\\ 0& A_{0101}& A_{0110}& 0\\ 0& A_{1001}& A_{1010}& 0\\ A_{1100}& 0& 0& A_{1111}\end{array}\right].$$ (20) ###### Proof: The construction of the twisting is similar as in lemma above. This time one has to consider also the singular value decomposition of the operator $`A_{0110}=WSW^{}`$. We can see now, that with any state $`\rho _{ABA^{}B^{}}`$, which has two qubit key part $`AB`$, we can associate a state obtained in the following way: 1. For state $`\rho _{ABA^{}B^{}}`$ find twisting $`U_{ps}`$, such, that (according to lemma 2) it changes upper-right element of $`AB`$ subsystem of $`\rho _{ABA^{}B^{}}`$ into $`A_{0011}`$. 2. Apply $`U_{ps}`$ to $`\rho _{ABA^{}B^{}}`$ obtaining $`\rho _{ABA^{}B^{}}^{}=U_{ps}\rho _{ABA^{}B^{}}U_{ps}^{}`$. 3. Trace out the shield ($`A^{}B^{}`$ subsystem) of state $`\rho _{ABA^{}B^{}}^{}`$ obtaining two-qubit state $$\rho _{AB}^{}=\mathrm{Tr}_{A^{}B^{}}\rho _{ABA^{}B^{}}^{}.$$ (21) This operation we will call privacy squeezing , or shortly p-squeezing , and the state $`\rho _{AB}^{}`$ which is the output of such operation on the state $`\rho _{ABA^{}B^{}}(𝒞^2𝒞^2𝒞^d𝒞^d^{})`$ the p-squeezed state of the state $`\rho _{ABA^{}B^{}}`$. Sometimes we shall use the term privacy squeezing in more informal sense, namely, with the twisting $`U_{ps}`$ which makes the key part close to maximally entangled state. Note, that the ccq state of p-squeezed state has no more secret correlations than that of the original state. This is because it emerges from the operation of twisting which preserves security in some sense, i.e. it does not change the ccq state which can be obtained from the original state. The next operation performed in definition of p-squeezed state is tracing out $`A^{}B^{}`$ part which means giving the $`A^{}B^{}`$ subsystem to Eve. Such operation can not increase security of the state in any possible sense. We will be interested in applying p-squeezing in the case, where the key part of the initial state was weakly entangled, or completely separable. Then the p-squeezing operation will make it entangled. We can say, that the operation of privacy squeezing pumps the entanglement of the state which is distributed along subsystems $`AA^{}BB^{}`$ into its key part $`AB`$. The entanglement once concentrated in the two qubit part, may be much more powerful than the one spread over the whole system. Further in the paper, we will see that from the bound entangled state, the operation of p-squeezing can produce approximately a maximally entangled state of two qubits. Then the analysis of how much key one can draw from the ccq state is much easier in the case of the p-squeezing state. ## III General form of states containing ideal key. In this section we will provide general form of the states $`\rho _{ABA^{}B^{}}`$ which have $``$ key, i.e. states such that the outcomes of measurement in basis $``$ are both perfectly correlated and perfectly secure. It turns out that this is precisely the class of private states. Hence the definitions 1 and 3 are equivalent. We have the following theorem: ###### Theorem 2 Any state $`\rho _{ABA^{}B^{}}`$ of a Hilbert space $`_A_A^{}_B_B^{}`$ with dimensions $`d_A=d_Bd`$, $`d_A^{}`$ and $`d_B^{}`$, has $``$-key if and only if it is of the form $$\rho _{ABA^{}B^{}}=\frac{1}{d}\underset{i,j=0}{\overset{d1}{}}\|e_if_ie_jf_j\|_{AB}U_i\sigma _{A^{}B^{}}U_j^{}$$ (22) where the state $`\sigma _{A^{}B^{}}`$ is an arbitrary state of subsystem $`A^{}B^{}`$, $`U_i`$’s are arbitrary unitary transformations and $`\{e_if_j\}=`$. We can rewrite the state (22) in the following, more appealing form $$\rho _{ABA^{}B^{}}=UP_{AB}^+\sigma _{A^{}B^{}}U^{},$$ (23) where $`P_{AB}^+`$ is a projector onto the maximally entangled state $`\psi _+=_i\frac{1}{\sqrt{d}}\|e_if_i`$, and $`U`$ is arbitrary twisting operation (5). Since the state $`P_{AB}^+`$ has many matrix elements vanishing, not all unitaries from definition of twisting are actually used here. In fact, unitaries $`U_i`$ in equation (22) are to be identified with unitaries $`U_{kk}`$ from equation (5). Note, that we can take $`\sigma _{A^{}B^{}}`$ to be ”classically correlated” in the sense that it is diagonal in some product basis. Indeed, twisting can change the state $`\sigma _{A^{}B^{}}`$ into any other state having the same eigenvalues (simply, twisting can incorporate a unitary transformation acting solely on $`A^{}B^{}`$). Thus we see that the states which have key, are closely connected with the maximally entangled state, which has been so far a ”symbol” of quantum security. As we shall see, the maximally entangled state may get twisted so much, that after measurement in many bases of the $`AB`$ part the outcomes will be correlated with Eve, which is not the case for the maximally entangled state itself. Still, however the basis $``$ will remain secure. Note that here we deal with perfect security. We will later discuss approximate security in Section V. ###### Proof: ($``$) This part of the proof is a consequence of the theorem 1. Namely a basic pdit (7) is obviously $``$-secure, because it has maximal correlations in this basis, and moreover it is a pure state, hence the one completely decoupled from Eve. More formally, it is evident that the ccq state of basic pdit is of the form (4). Now we can apply theorem 1, which says that after twisting the ccq state is unchanged. Hence any state of the form (22) has also $``$ key. ###### Proof: ($``$) In this part we assume, that the state $`\rho _{ABA^{}B^{}}`$ has $``$-key i.e. that after measurement on it’s $`AB`$ part, one gets a perfectly correlated state (between Alice and Bob) that is uncorrelated with Eve: $$\left(\underset{i=0}{\overset{d1}{}}\frac{1}{d}\|e_if_ie_if_i\|_{AB}\right)\rho _E.$$ (24) Let us consider general pure state for which dimensions of $`A,B`$ are $`d`$, dimensions of $`A^{},B^{}`$ are $`d_A^{},d_B^{}`$ respectively, and dimension of subsystem $`E`$ is the smallest one which allows for the whole state being a pure one. $$\|\psi =\|\psi _{ABA^{}B^{}E}=\underset{ijklm}{}a_{ijklm}\|e_if_jklm.$$ (25) one can rewrite it as $$\|\psi =\underset{ij}{}\|e_if_j_{AB}\|\stackrel{~}{\psi }^{(ij)}_{A^{}B^{}E}.$$ (26) with $`\|\stackrel{~}{\psi }^{(ij)}_{A^{}B^{}E}=_{klm}a_{ijklm}\|klm`$. It is easy to see that the scalar product $`\stackrel{~}{\psi }^{(ij)}\|\stackrel{~}{\psi }^{(ij)}`$ equals the probability of obtaining the state $`\|e_if_je_if_j\|_{AB}`$ on the system $`AB`$ after measurement in basis $``$. Now, since the subsystem $`\rho _{ABE}`$ (after measurement in $``$ on AB) must be maximally correlated, the vectors $`\|\stackrel{~}{\psi }^{(ij)}`$ should satisfy $`\stackrel{~}{\psi }^{(ij)}\|\stackrel{~}{\psi }^{(ij)}=\frac{1}{d}\delta _{ij}`$. We can normalize these states (in case $`i=j`$) to have: $$\|\psi ^{(ii)}:=\frac{\|\stackrel{~}{\psi }^{(ii)}}{\sqrt{\stackrel{~}{\psi }^{(ii)}\|\stackrel{~}{\psi }^{(ii)}}}=\sqrt{d}\|\stackrel{~}{\psi }^{(ii)}$$ (27) so that the total state has a form: $$\|\psi =\underset{i=0}{\overset{d1}{}}\frac{1}{\sqrt{d}}\|e_if_i_{AB}\|\psi ^{(ii)}_{A^{}B^{}E}.$$ (28) ”Cryptographical” interpretation of this state is the following: if Alice and Bob gets $`i`$th result, then Eve gets subsystem $`\rho _i^E`$ of a state $`\|\psi ^{(ii)}_{A^{}B^{}E}`$. Indeed, the ccq state is then of the form $$\rho _{ccq}=\underset{i=0}{\overset{d1}{}}\frac{1}{d}\|e_if_i_{AB}e_if_i\|\rho _i^E,$$ (29) with $`\rho _i=\mathrm{Tr}_{A^{}B^{}}(\|\psi ^{(ii)}_{A^{}B^{}E}\psi ^{(ii)}\|)`$. Now the condition (24) implies that, $`\rho _i^E`$ should be all equal to each other. In particular, it follows that rank of Eve’s total density matrix is no greater than dimension of $`A^{}B^{}`$ system, hence we can assume that $`d_E=d_A^{}d_B^{}=d^{}`$. It is convenient to rewrite this pure state in a form $$\|\psi ^{(ii)}_{A^{}B^{}E}=\underset{k=0}{\overset{d^{}1}{}}\|k_{A^{}B^{}}X_i\|k_E,$$ (30) where $`\{\|k\}`$ is standard basis of $`A^{}B^{}`$ and of $`E`$ system, $`X_i`$ is $`d_E\times d_E`$ matrix that fully represents this state. It is easy to check, that $`\rho _i^E=X_iX_i^{}`$. Consider now singular value decomposition of $`X_i`$ given by $`V_i\sqrt{\rho _i}U_i^{}`$ where $`\rho _i`$ is now diagonal in basis $`\{\|k\}`$. One then gets that $`\rho _i^E=V_i\rho _iV_i^{}`$. The state (30) may be rewritten $$\|\psi ^{(ii)}_{A^{}B^{}E}=\underset{k}{}X_i^T\|k_{A^{}B^{}}\|k_E,$$ (31) where $`T`$ is transposition in basis $`\{\|k\}`$. Now it is easy to check, that subsystem $`A^{}B^{}`$ of $`\|\psi ^{(ii)}_{A^{}B^{}E}`$ is in state $`X_i^T(X_i^T)^{}`$, so that the whole state $`\rho _{ABA^{}B^{}}`$ is the following: $$\rho _{ABA^{}B^{}}=\frac{1}{d}\underset{i,j=0}{\overset{d1}{}}\|e_if_ie_jf_j\|_{AB}X_i^T(X_j^{})^T.$$ (32) We can express this state using states accessible to Eve, namely $`\rho _j^E`$: $`\rho _{ABA^{}B^{}}={\displaystyle \frac{1}{d}}{\displaystyle \underset{i,j=0}{\overset{d1}{}}}\|e_if_ie_jf_j\|_{AB}`$ $`(U_i^{}V_i^T)\underset{=\sqrt{\rho _i^E}^T}{\underset{}{V_i^{}\sqrt{\rho _i}^TV_i^T}}\underset{=\sqrt{\rho _j^E}^T}{\underset{}{V_j^{}\sqrt{\rho _j}^TV_j^T}}(V_j^{}U_j^T).`$ (33) (For expressing state in terms of Eve’s states in more general case, see Section VI). Denoting by $`W_i`$ the unitary transformation $`U_i^{}V_i^T`$ one gets: $$\rho _{ABA^{}B^{}}=\frac{1}{d}\underset{i,j=0}{\overset{d1}{}}\|e_if_ie_jf_j\|_{AB}W_i\sqrt{\rho _i^E}^T.\sqrt{\rho _j^E}^TW_j^{}.$$ However, as mentioned above, Eve’s density matrices are equal to each other, i.e. $`\rho _i^E=\rho _j^E`$ for all $`i,j`$. We then obtain $$\rho _{ABA^{}B^{}}=\frac{1}{d}\underset{i,j=0}{\overset{d1}{}}\|e_if_ie_jf_j\|_{AB}W_i\rho W_{j}^{}{}_{A^{}B^{}}{}^{}.$$ (34) This completes the proof of theorem 2. ## IV Pdits and their properties In this section we will present various forms of pdits and pbits. We will first write the pbit in matrix form according to its original definition. We can write it in block form $$\gamma _{ABA^{}B^{}}^{(2)}=\frac{1}{2}\left[\begin{array}{cccc}U_0\sigma _{A^{}B^{}}U_0^{}& 0& 0& U_0\sigma _{A^{}B^{}}U_1^{}\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ U_1\sigma _{A^{}B^{}}U_0^{}& 0& 0& U_1\sigma _{A^{}B^{}}U_1^{}\end{array}\right],$$ (35) where $`\sigma _{A^{}B^{}}`$ is arbitrary state on $`A^{}B^{}`$ subsystem, and $`U_0`$ and $`U_1`$ are arbitrary unitary transformations which act on $`A^{}B^{}`$. ”Generalized EPR form” of pdit. Since by the theorem of the previous section pdits are the only states that contain $``$-key, they could be called generalized EPR states (maximally entangled state). We have already seen that they are ”twisted EPR states”. One can notice an even closer connection. Namely, a pdit can be viewed as an EPR states with operator amplitudes. Indeed, one can rewrite equation (32) in a more appealing form $$\gamma _{ABA^{}B^{}}^{(d)}=\mathrm{\Psi }\mathrm{\Psi }^{},$$ (36) with $$\mathrm{\Psi }=\frac{1}{\sqrt{d}}\underset{i=0}{\overset{d1}{}}Y_i^{A^{}B^{}}\|e_if_i_{AB}.$$ (37) We have written here (unlike in the rest the of paper) first the $`A^{}B^{}`$ system and then the $`AB`$ one, so that this form of pdit would recall a form of pure state. Thus instead of $`c`$-numbers the amplitudes are now $`q`$-numbers, so that states which have key are ”second quantized EPR states”. In the case of pbits, the matrix form is the following: $$\gamma _{ABA^{}B^{}}^{(2)}=\frac{1}{2}\left[\begin{array}{cccc}Y_0Y_0^{}& 0& 0& Y_0Y_1^{}\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ Y_1Y_0^{}& 0& 0& Y_1Y_1^{}\end{array}\right].$$ (38) Let us consider the polar decomposition of operators $`Y_i`$. From definition of pdit it follows that the only constraint on these operators is the following $$_iY_i=U_i\sqrt{\rho },$$ (39) where $`U_i`$ is unitary transformation and $`\rho `$ is a normalized state as so is the $`\sigma _{A^{}B^{}}`$ state in form (35). This reflects the fact, that similarly like maximally entangled state which produces correlated outputs has amplitudes with probably different phases, but the same moduli, the ”maximally private” state can have ”operator amplitudes” which differ by $`U_i`$ ( a counterpart of the phase) but have the same $`\sqrt{\rho }`$ in polar decomposition (which is a counterpart of the modulus). There is yet another similarity to EPR states, namely the norm of upper-right block $`Y_0Y_1^{}`$ is equal to $`\frac{1}{2}`$, like the modulus of the coherence of the EPR state. $`X`$-form” of pbit. In special case of pbits one can have representation by just one normalized operator: $$\gamma _{ABA^{}B^{}}^{(2)}=\frac{1}{2}\left[\begin{array}{cccc}\sqrt{XX^{}}& 0& 0& X\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ X^{}& 0& 0& \sqrt{X^{}X}\end{array}\right],$$ (40) for any operator $`X`$ satisfying $`X=1`$. Justification of equivalence of this form and standard form is the following. Consider singular value decomposition of $`X`$ $`X=U\sigma W`$ with $`U`$ and $`W`$ unitary transformations and $`\sigma `$ being diagonal, positive matrix. Since $`X`$ has trace norm 1, the same is for $`\sigma `$, therefore it can be viewed as $`X=U\rho W`$ with $`\rho `$ being a legitimate state. Identifying $`U_0=U`$ and $`U_1=W^{}`$ we obtain standard form. It is important, that in nontrivial cases $`X`$ should be a non-positive operator. Otherwise the pbit is equal to basic pbit. Indeed, if it is positive, then since its trace norm is 1, it is itself a legitimate state, call it $`\rho `$. Then $`\sqrt{XX^{}}=\sqrt{X^{}X}=\rho `$, so that $$\rho _{ABA^{}B^{}}=\frac{1}{2}\underset{i,j=0}{\overset{1}{}}\|iijj\|\rho =\|\psi _+\psi _+\|\rho .$$ which is the basic pbit (7). Note, that in higher dimension to have the $`X`$-form we need more than one operator, and the operators depend on each other, which is not as simple representation as in the case of pbit. For example in $`d=3`$ case we have: $`{\displaystyle \frac{1}{3}}\left[\begin{array}{ccccccccc}\sqrt{XX^{}}& 0& 0& 0& X& 0& 0& 0& XY\\ 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0\\ X^{}& 0& 0& 0& \sqrt{X^{}X}& 0& 0& 0& Y\\ 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0\\ (XY)^{}& 0& 0& 0& Y^{}& 0& 0& 0& \sqrt{Y^{}Y}\end{array}\right],`$ (50) where the operators $`X`$ and $`Y`$ satisfy: $`X=1`$ and $`X=WY^{}`$ for arbitrary unitary transformation $`W`$. ”Flags form”: special case of $`X`$-form. If the operator $`X`$ which represents pbit in its $`X`$-form is additionally hermitian, any such pbit can be seen as a mixture of basic pbit and a variation of basic pbit which has EPR states with different phase: $$\gamma _{ABA^{}B^{}}^{(2)}=p\|\psi _+\psi _+\|\rho _{A^{}B^{}}^++(1p)\|\psi _{}\psi _{}\|\rho _{A^{}B^{}}^{},$$ (51) where $`\|\psi _\pm =\frac{1}{\sqrt{2}}(\|00\pm \|11)`$. Derivation of this form is straightforward, if we consider decomposition of X into positive and negative part : $$X=X_+X_{},$$ (52) where $`X_+`$ and $`X_{}`$ are by definition orthogonal, and positive. Thus denoting $`p=\mathrm{Tr}X_+`$, together with assumption of $`X`$-form that $`X=\mathrm{Tr}\|X\|=1`$, we can rewrite $`X`$ as $$X=p\rho _+(1p)\rho _{},$$ (53) where $`\rho _\pm `$ are normalized positive and negative parts of $`X`$. Moreover, since the states $`\rho _+`$ and $`\rho _{}`$ are orthogonal: $`\mathrm{Tr}\rho _{}\rho _+=0`$, we obtain the form (51). ### IV-A Private bits - examples We will give now two examples of private bits, and study its entanglement distillation properties. Examples of pbit 1. Let us consider state $`\gamma ^VB(𝒞^2𝒞^2𝒞^d𝒞^d)`$ of the following form: $$\gamma ^V=\frac{1}{2}\left[\begin{array}{cccc}\frac{I}{d^2}& 0& 0& \frac{V}{d^2}\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ \frac{V}{d^2}& 0& 0& \frac{I}{d^2}\end{array}\right],$$ (54) where $`V`$ is the swap operator which reads: $`V=_{i=0}^{d1}\|ijji\|`$. If we consider positive and negative part of $`V`$, which are symmetric and antisymmetric subspace, it is easy to see, that $$\gamma ^V=p\|\psi _+\psi _+\|\rho _s+(1p)\|\psi _{}\psi _{}\|\rho _a$$ (55) where $$\rho _s=\frac{2}{d^2+d}P_{sym}\rho _a=\frac{2}{d^2d}P_{asym}$$ (56) are symmetric and antisymmetric Werner states, and $`p=\frac{1}{2}(1+\frac{1}{d}`$). Thus we have obtained, that it is also a pbit with natural ”flags form”, with flags being (orthogonal) Werner states . 2. The second example is the state known as ”flower state”, which was shown to lock entanglement cost. We have that $`\rho _{flower}B(𝒞^2𝒞^2𝒞^{d^2}𝒞^{d^2})`$ is of the form: $$\gamma _{flower}=\frac{1}{2}\left[\begin{array}{cccc}\sigma & 0& 0& \frac{1}{d}U^T\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ \frac{1}{d}U^{}& 0& 0& \sigma \end{array}\right],$$ (57) where $`\sigma `$ is classical maximally correlated state: $`\sigma =_{i=0}^{d1}\frac{1}{d}\|iiii\|`$, and $`U`$ is the embedding of unitary transformation $`W=_{i,j=0}^{d1}w_{ij}\|ij\|=H^{\mathrm{log}d}`$ with $`H`$ being Hadamard transform in the following way: $$U=\underset{i,j=0}{\overset{d1}{}}w_{ij}\|iijj\|.$$ We can check now, that this state is pbit with $`X`$-form. In this case $`X=U^T`$. To see this consider unitary transformation $`S:=U^{}+_{ij}\|ijij\|`$. Composing $`S`$ with $`U^T`$ does not change the norm, which is unitarily invariant, so that $$\frac{1}{d}U^T=\frac{1}{d}U^TS=\frac{1}{d}\underset{i=0}{\overset{d1}{}}\|iiii\|=1.$$ (58) Thus we see, that $`X=1`$. We have also $`\sqrt{XX^{}}=\sigma `$: $$\sqrt{\frac{1}{d^2}U^TU^{}}=[\frac{1}{d^2}\underset{i=0}{\overset{d1}{}}\|iiii\|]^{\frac{1}{2}}=\sigma .$$ (59) We will show now, that in case of $`\gamma ^V`$ given in eq. (54), the distillable entanglement $`E_D`$ is strictly smaller then the amount of secure key $`K_D`$ gained from these states. The formal definition of $`K_D`$ is given in Section VIII. Here it is enough to base only on its intuitive properties. Namely, any pdit by its very definition has $`K_D`$ at least equal to $`\mathrm{log}d`$ of key, which can be obtained by measuring its key part. To show the gap between distillable entanglement and distillable key we will compute the value of another measure of entanglement: log-negativity $`E_N(\rho )`$ (see ) of the state, which is an upper bound on distillable entanglement . To this end consider the following lemma. ###### Lemma 3 For any pbit in $`X`$-form, if $`\sqrt{XX^{}}`$ and $`\sqrt{X^{}X}`$ are PPT, the log negativity of the pbit in $`X`$-form reads $`E_N=\mathrm{log}(1+X^\mathrm{\Gamma })`$, where $`\mathrm{\Gamma }`$ is transposition performed on the system $`B^{}`$. The proof of this lemma is given in Sec. XV-E of Appendix. Using this lemma, one can check the negativity of the state $`\gamma ^V`$. We have in this case $`X=\frac{V}{d^2}`$, with $`d2`$. Since $`V^\mathrm{\Gamma }=dP_+`$, we obtain $`E_N(\gamma ^V)=\mathrm{log}(1+\frac{1}{d})`$. It implies: $$E_D(\gamma ^V)E_N(\gamma ^V)=\mathrm{log}(1+\frac{1}{d})<1K_D(\gamma ^V),$$ (60) which demonstrates a desired gap between distillable key and distillable entanglement: $$E_D(\gamma ^V)<K_D(\gamma ^V).$$ (61) ### IV-B Relative entropy of entanglement and pdits In this section we will consider the entanglement contents of the pbit in terms of a measure of entanglement called relative entropy of entanglement, defined as follows: $$E_r(\rho )=\underset{\sigma _{sep}SEP}{inf}S(\rho \|\sigma _{sep})$$ (62) where $`S(\rho \|\|\sigma )=S(\rho )\mathrm{Tr}\rho \mathrm{log}\sigma `$ is the relative entropy, and $`SEP`$ is the set of separable states. In the Section XI we will show, that for any state, the relative entropy of entanglement is an upper bound on the key rate, that can be obtained from the state (for generalizations of this result to a wide class of entanglement monotones see ). It is then easy to see, that for any pbit $`\gamma `$, $`E_r(\gamma )`$ is greater than $`\mathrm{log}d`$ since $`K_D(\gamma )\mathrm{log}d`$ by definition of pdits. The question we address here, is the upper bound on the relative entropy of the pdit. We relate its value to the states which appear on the shield of the pdits, when Alice and Bob get key by measuring the key part of the pdit. The theorem below states it formally. ###### Theorem 3 For any pdit $`\gamma _{ABA^{}B^{}}(𝒞^d𝒞^d𝒞^{d_A^{}}𝒞^{d_B^{}})`$, which is secure in standard basis, let $`\rho _{A^{}B^{}}^{(i)}`$ denote states which appears on shield of the pbit, after obtaining outcome $`ii`$ in measurement performed in standard basis on its key part. Then we have $$E_r(\gamma _{ABA^{}B^{}})\mathrm{log}d+\frac{1}{d}\underset{i=0}{\overset{d1}{}}E_r(\rho _{A^{}B^{}}^{(i)})$$ (63) where $`AB`$ denotes key and $`A^{}B^{}`$ shieldpart of the pdit. ###### Proof: One can view the quantity $`\frac{1}{d}_{i=0}^{d1}E_r(\rho _{A^{}B^{}}^{(i)})`$ as the relative entropy of $`\gamma _{ABA^{}B^{}}`$ dephased on $`AB`$ in computational basis . In case $`d=2^k`$, it can be easily done with applying unitary $`U_i`$ \- random sequence of $`\sigma _z`$ and $`I`$ unitary transformations. In general case one can use the so called Weyl unitary operators (see e.g. ). Such an implementation of dephasing uses $`\mathrm{log}d`$ bits of randomness. Following the proof of non-lockability of relative entropy of entanglement (see also ), we can write $$E_r(\gamma _{ABA^{}B^{}})E_r(\underset{i}{}p_i\sigma _i)\mathrm{log}d$$ (64) where $`\sigma _i=U_iI_{A^{}B^{}}\gamma _{ABA^{}B^{}}U_i^{}I_{A^{}B^{}}`$ and $`p_i=\frac{1}{d}`$. As we have observed above, the relative entropy of dephased state $`_ip_i\sigma _i`$ equals $`\frac{1}{d}_{i=0}^{d1}E_r(\rho _{A^{}B^{}}^{(i)})`$ which ends the proof. The above theorem is valid also for regularized relative entropy, defined as $$E_r(\rho )^{\mathrm{}}=\underset{n\mathrm{}}{lim}\frac{1}{n}E_r(\rho ^n).$$ (65) ###### Theorem 4 Under the assumptions of theorem 3 there holds: $$E_r^{\mathrm{}}(\gamma _{ABA^{}B^{}})\mathrm{log}d+\frac{1}{d}\underset{i=0}{\overset{d1}{}}E_r^{\mathrm{}}(\rho _{A^{}B^{}}^i),$$ (66) For the proof of this theorem, as rather technical, we refer the reader to Appendix XV-F. ### IV-C Irreducible pbit - a unit of privacy In Section III we have characterized states which contain ideal key i.e. pdits. A pdit has an $`AB`$ subsystem called here the key part. $`\mathrm{log}d`$ bits of key can be obtained from such a pdit by a complete measurement in some basis performed on this key part of pdit . However, as it follows from the characterization given in theorem 2, pdits have also the $`A^{}B^{}`$ subsystem, called here the shield. This part can also serve as a source of key. Indeed there are plenty of such pdits that contain more than $`\mathrm{log}d`$ key, due to their shield. Therefore not every pdit can serve as a of unit of privacy and we need the following definition: ###### Definition 5 Any pdit $`\gamma `$ (with $`d`$-dimensional key part) for which $`K_D(\gamma )=\mathrm{log}d`$ is called irreducible. This definition distinguishes those $`pdits`$ for which measuring their key part is the optimal protocol for drawing key. They are called irreducible in opposite to those, which can be reduced by distillation protocol to some other pdits which have more than $`\mathrm{log}d`$ of key. Irreducible pdits are by definition units of privacy (although they are not generally interconvertible). Determining the class of irreducible pdits is potentially a difficult task, as it leads to optimisation over protocols of key distillation. However we are able to show a subclass of pdits, which are irreducible. To this end we use a result, which is proven in Section XI, namely that the relative entropy of entanglement is an upper bound on distillable key. Having this we can state the following proposition: ###### Proposition 1 Any pdit $`\gamma `$, with $`E_r(\gamma )=\mathrm{log}d`$, is irreducible. ###### Proof: By definition of $`pdit`$ we have $`K_D(\gamma )\mathrm{log}d`$ and by theorem 169 from Section XI we have $`K_D(\gamma )E_r(\gamma )`$ which is in turn less than $`\mathrm{log}d`$ by assumption, and the assertion follows. We can provide now a class of pdits which have $`E_r=\mathrm{log}d`$ and by the above proposition are irreducible. These are pdits which have separable states that appear on shield conditionally on outcomes of complete measurement on key part part in the computational basis. ###### Proposition 2 For any pdit $`\gamma _{ABA^{}B^{}}(𝒞^d𝒞^d𝒞^{d_A^{}}𝒞^{d_B^{}})`$, which is secure in standard basis, if $`\rho _{A^{}B^{}}^{(i)}`$ denote states which appear on shield of the pbit, after obtaining outcome $`\|i\|i`$ in measurement performed in standard basis on its key part are separable states, then pdit $`\gamma _{ABA^{}B^{}}`$ is irreducible. ###### Proof: Due to bound on relative entropy of pdit given in theorem 3 we have that $`E_r(\gamma )`$ is less then or equal to $`\mathrm{log}d`$ since conditional states $`\rho _{A^{}B^{}}^{(i)}`$ are separable and hence have relative entropy of entanglement equal to zero. $`E_r(\gamma )`$ is also not less then $`\mathrm{log}d`$, since it is greater than the amount of distillable key, which ends the proof. Note, that examples (54), (57) given in Section IV-A fulfill the assumptions of this theorem, and are therefore irreducible pbits. They are also the first known non trivial states (different than pure state) for which the amount of distillable key has been calculated. Using the bound of relative entropy on distillable key, one can also show, that the class of maximally correlated states has $`K_D=E_D=E_r`$, since for the latter $`E_D=E_r`$. ## V Approximate pbits We present here a special property of states which are close to pbit. We have already seen, that pbits have similar properties to the maximally entangled EPR states. In particular, the norm of the upper-right block in standard form as well as in $`X`$-form of pbit is equal to $`\frac{1}{2}`$. We will show here, that for general states the norm of that block tells how close the state is to a pbit: any state which is close in trace norm to pbit must have the norm of this block close to $`\frac{1}{2}`$, and vice versa. We will need the following lemma that relates the value of coherence to the distance from the maximally entangled state for two qubit states. ###### Lemma 4 For any bipartite state $`\rho _{AB}B(𝒞^2𝒞^2)`$ expressed on the form $`\rho _{AB}=_{ijkl=0}^1a_{ijkl}\|ijkl\|`$ we have: $$\mathrm{Tr}\rho _{AB}P_+1ϵRe(a_{0011})>\frac{1}{2}ϵ$$ (67) and $$Re(a_{0011})>\frac{1}{2}ϵ\mathrm{Tr}\rho _{AB}P_+12ϵ$$ (68) ###### Proof: For the proof of this lemma, see Appendix XV-G. We can prove now that approximate pbits have norm of an appropriate block close to $`\frac{1}{2}`$. ###### Proposition 3 If the state $`\sigma _{ABA^{}B^{}}B(𝒞^2𝒞^2𝒞^d𝒞^d^{})`$ written in the form $`\sigma _{ABA^{}B^{}}=_{ijkl=0}^1\|ijkl\|A_{ijkl}`$ fulfills $$\sigma _{ABA^{}B^{}}\gamma _{ABA^{}B^{}}ϵ$$ (69) for some pbit $`\gamma `$, then for $`0<ϵ<1`$ there holds $`A_{0011}\frac{1}{2}ϵ`$. ###### Proof: The pbit $`\gamma `$ is a twisted EPR state, which means that there exists twisting $`U`$ which applied to basic pbit $`P_+\rho `$ gives $`\gamma `$. We apply this $`U`$ to both states $`\sigma _{ABA^{}B^{}}`$ and $`\gamma `$ and trace out the $`A^{}B^{}`$ subsystem of both of them. Since these operations can not increase the norm distance between these states, so that we have for $`\sigma _{AB}=\mathrm{Tr}_{A^{}B^{}}U\sigma _{ABA^{}B^{}}U^{}`$ $$\sigma _{AB}P_+ϵ.$$ (70) It implies, by equivalence of norm and fidelity (1) that $$F(\sigma _{AB},P_+)1\frac{1}{2}ϵ.$$ (71) We have also that $`F(\sigma _{AB},P_+)^2=\mathrm{Tr}\sigma _{AB}P_+`$ so that $$\mathrm{Tr}\sigma _{AB}P_+>1ϵ$$ (72) for $`ϵ<1`$. Now by lemma (68) this yields $`\|a_{0011}\|Re(a_{0011})\frac{1}{2}ϵ`$, where $`a_{0011}`$ is coherence of the state $`\rho _{AB}=_{ijkl=0}^1a_{ijkl}\|ijkl\|`$. However, we have $$\|a_{0011}\|=\|\mathrm{Tr}U_{00}A_{0011}U_{11}^{}\|$$ (73) where $`U_{00}`$ and $`U_{11}`$ come from twisting, that we have applied. Using now the fact that $`A=sup_U\mathrm{Tr}AU`$, where supremum is taken over unitary transformations we get $$A_{0011}\|a_{0011}\|1ϵ.$$ (74) This ends the proof. Now we will formulate and prove the converse statement, saying that when the norm of the right upper block is close to $`1/2`$, then the state is close to some pbit. ###### Proposition 4 If the state $`\sigma _{ABA^{}B^{}}B(𝒞^2𝒞^2𝒞^d𝒞^d^{})`$ with a form $`\sigma _{ABA^{}B^{}}=_{ijkl=0}^1\|ijkl\|A_{ijkl}`$ fulfills $`A_{0011}>\frac{1}{2}ϵ`$ then for $`0<ϵ<\frac{1}{8}`$ there exists pbit $`\gamma `$ such, that $$\sigma _{ABA^{}B^{}}\gamma _{ABA^{}B^{}}\delta (ϵ)$$ (75) with $`\delta (ϵ)`$ vanishing, when $`ϵ`$ approaches zero. More specifically, $$\delta (ϵ)=2\sqrt{8\sqrt{2ϵ}+h(2\sqrt{2ϵ})}+2\sqrt{2ϵ}$$ (76) with $`h(.)`$ being the binary entropy function $`h(x)=x\mathrm{log}x(1x)\mathrm{log}(1x)`$. ###### Proof: In this proof by $`\rho _X`$ we denote respective reduced density matrix of the state $`\rho _{ABA^{}B^{}}`$. Let $`\rho _{AB}`$ be the privacy-squeezed state of the state $`\sigma _{ABA^{}B^{}}`$ i.e. $`\rho _{AB}=Tr_{A^{}B^{}}\rho _{ABA^{}B^{}}`$ where $`\rho _{ABA^{}B^{}}=U_{ps}\sigma _{ABA^{}B^{}}U_{ps}^{}`$ for certain twisting $`U_{ps}`$. By definition the entry $`a_{0011}`$ of $`\rho _{AB}`$ is equal to $`A_{0011}`$. By assumption we have, $`a_{0011}=A_{0011}>\frac{1}{2}ϵ`$. By lemma 68 (equation (68)) we have that $$\mathrm{Tr}\rho _{AB}P_+>12ϵ.$$ (77) We have then $$F(\rho _{AB},P_+)^2=\mathrm{Tr}\rho _{AB}P_+$$ (78) which, by equivalence of norm and fidelity (11) gives $$\rho _{AB}P_+2\sqrt{2ϵ}.$$ (79) Let us now consider the state $`\rho _{ABA^{}B^{}}=U_{ps}\sigma _{ABA^{}B^{}}U_{ps}^{}`$ and its purification to Eve’s subsystem $`\psi _{ABA^{}B^{}E}`$ so that we have: $$\rho _{AB}=\mathrm{Tr}_{A^{}B^{}E}(\psi _{ABA^{}B^{}E})$$ (80) By the Fannes inequality (see eq. (12) in Sec. II-A) we have that $$S(\rho _{AB})=S(\rho _{A^{}B^{}E})4\sqrt{2ϵ}\mathrm{log}d_{AB}+h(2\sqrt{2ϵ}).$$ (81) ¿From this we will get that $`\psi _{ABA^{}B^{}E}\rho _{AB}\rho _{A^{}B^{}E}`$ is of order of $`ϵ`$. We prove this as follows. Since norm distance is bounded by relative entropy as follows $$\frac{1}{2}\rho _1\rho _2^2S(\rho _1\|\rho _2),$$ (82) one gets: $$\psi _{ABA^{}B^{}E}\rho _{AB}\rho _{A^{}B^{}E}\sqrt{2S(\psi _{ABA^{}B^{}E}\|\|\rho _{AB}\rho _{A^{}B^{}E})}.$$ The relative entropy distance of the state to it’s subsystems is equal to quantum mutual information $$I(\psi _{AB\|A^{}B^{}E})=S(\rho _{AB})+S(\rho _{A^{}B^{}E})S(\psi _{ABA^{}B^{}E}),$$ (83) (We will henceforth use shorthand notation $`I(X:Y)`$, $`S(X)`$). which gives $$I(AB:A^{}B^{}E)=2S(AB)2(4\sqrt{2ϵ}\mathrm{log}d_{AB}+h(2\sqrt{2ϵ})).$$ (84) where last inequality comes from Eq. (81). Coming back to inequality (V) we have that $$\psi _{ABA^{}B^{}E}\rho _{AB}\rho _{A^{}B^{}E}\sqrt{2I(AB:A^{}B^{}E)}2\sqrt{4\sqrt{2ϵ}\mathrm{log}d_{AB}+h(2\sqrt{2ϵ})}$$ (85) If we trace out the subsystem $`E`$ the inequality is preserved: $$\rho _{ABA^{}B^{}}\rho _{AB}\rho _{A^{}B^{}}2\sqrt{8\sqrt{ϵ}+h(2\sqrt{ϵ})},$$ (86) where we have put $`d_{AB}=4`$, as we deal with pbits. Now by triangle inequality one has: $$\rho _{ABA^{}B^{}}P_+\rho _{A^{}B^{}}\rho _{ABA^{}B^{}}\rho _{AB}\rho _{A^{}B^{}}+\rho _{AB}\rho _{A^{}B^{}}P_+\rho _{A^{}B^{}}.$$ (87) We can apply now the bounds (79) and (86) to the above inequality obtaining $$\rho _{ABA^{}B^{}}P_+\rho _{A^{}B^{}}2\sqrt{8\sqrt{2ϵ}+h(2\sqrt{2ϵ})}+2\sqrt{2ϵ}.$$ (88) Let us now apply the twisting $`U_{ps}^{}`$ (transformation which is inverse to twisting $`U_{ps}`$) to both states on left-hand-side of the above inequality. Since $`\rho _{ABA^{}B^{}}`$ is defined as $`U_{ps}\sigma _{ABA^{}B^{}}U_{ps}^{}`$ we get that: $$\sigma _{ABA^{}B^{}}U_{ps}^{}P_+\rho _{A^{}B^{}}U_{ps}2\sqrt{8\sqrt{2ϵ}+h(2\sqrt{2ϵ})}+2\sqrt{2ϵ},$$ (89) i.e. our state is close to pbit $`\gamma =U_{ps}^{}P_+\rho _{A^{}B^{}}U_{ps}`$. Then the theorem follows with $`\delta (ϵ)=2\sqrt{8\sqrt{2ϵ}+h(2\sqrt{2ϵ})}+2\sqrt{2ϵ}`$. ###### Remark 2 The above propositions establish the norm of upper-right block of matrix (written in computational basis according to ABA’B’ order of subsystems), as a parameter that measures closeness to pbit, and in this sense it measures security of the bit obtained from the key part. The state of form (19) is close to a pbit if and only if the norm of this block is close to $`\frac{1}{2}`$. This is the property of approximate pbits, however it seems not to have an analogue for approximate pdits with $`d3`$. ## VI Expressing Alice and Bob states in terms of Eve’s states In this section we will express the state $`\rho _{ABA^{}B^{}}`$ in such a way that one explicitly sees Eve’s states in it. We will then interpret the results of the previous sections in terms of such a representation. In particular, we will see that the norm of the upper-right block not only measures closeness to pbit, but it also measures the security of the bit from the key part directly, in terms of fidelity between corresponding Eve’s states. ### VI-A The case without shield. ”Abelian” twisting. Consider first the easier case of a state without shield i.e. $$\rho _{AB}=\underset{iji^{}j^{}}{}\rho _{iji^{}j^{}}\|iji^{}j^{}\|.$$ (90) Purification of this state is of the following form $$\psi _{ABE}=\underset{ij}{}\sqrt{p_{ij}}\|ij_{AB}\|\psi _E^{ij}$$ (91) where $`p_{ij}=\rho _{ijij}`$. We see, that when Alice and Bob measure the state in basis $`\|ij`$, Eve’s states corresponding to outcomes $`ij`$ are $`\psi _{ij}`$, and they occur with probabilities $`p_{ij}`$. Performing partial trace over Eve’s system, one obtains $$\rho _{AB}=\underset{iji^{}j^{}}{}\sqrt{p_{ij}p_{i^{}j^{}}}\psi _E^{i^{}j^{}}\|\psi _E^{ij}\|iji^{}j^{}\|$$ (92) Thus the matrix elements of $`\rho _{AB}`$ are inner products of Eve’s states. If we have all inner products between set of states, we have complete knowledge about the set, up to a total unitary rotation, which is irrelevant for security issues (since Eve can perform this herself). Thus density matrix $`\rho _{AB}`$ can be represented in such a way that all properties of Eve’s states are explicitly displayed. Moreover, moduli of matrix elements are related to fidelity between Eve’s states: $$\|\rho _{iji^{}j^{}}\|=\sqrt{p_{ij}p_{i^{}j^{}}}F(\psi _E^{ij},\psi _E^{i^{}j^{}}).$$ (93) #### Two-qubit case. For example, for two qubits, the density matrix looks as follows (we have not shown all elements) $$\rho _{AB}=\left[\begin{array}{cccc}p_{00}& \times & \times & \sqrt{p_{00}p_{11}}\psi _E^{11}\|\psi _E^{00}\\ \times & p_{01}& \times & \times \\ \times & \times & p_{10}& \times \\ \times & \times & \times & p_{11}\end{array}\right]$$ (94) Let us now consider the conditions for having one bit of perfect key obtained from the measurement in the two qubit case. They are as follows: (i) $`p_{00}=p_{11}=1/2`$ and (ii) $`\psi _{00}=\psi _{11}`$ up to a phase factor. The latter condition is equivalent to $`F(\psi _{00},\psi _{11})=1`$ (we have dropped here the index $`E`$). The two conditions can be represented by a single condition: $$\sqrt{p_{00}p_{11}}F(\psi _{00},\psi _{11})=\frac{1}{2}$$ (95) However, we know from (93) that this means that upper-right matrix element of $`\rho _{AB}`$ should satisfy $`\|\rho _{0011}\|=1/2`$. Consider now approximate bit of key, so that the conditions are satisfied up to some accuracy. Again we can combine them into single condition $$\sqrt{p_{00}p_{11}}F(\psi _{00}^E,\psi _{11}^E)>\frac{1}{2}ϵ$$ (96) This translates into $$\|\rho _{0011}\|\frac{1}{2}ϵ,$$ (97) which gives: $$F^2\frac{1}{2}(1+2\|\rho _{0011}\|),$$ (98) where $`\psi _{ME}=\frac{1}{\sqrt{2}}(e^{i\varphi _{00}}\|00+e^{i\varphi _{11}}\|11)`$, and $`F^2=\psi _{ME}\|\rho \|\psi _{ME}`$. Thus, in particular for $`F=1`$, we must have $`\|\rho _{0011}\|=1/2`$. The change of phases can be viewed as a unitary operation, where phases are controlled by the basis $`\|ij`$: $$U=\underset{ij}{}\|ijij\|e^{i\varphi _{ij}}$$ (99) Since we have $`\psi _{\mathrm{max}}=U\|\psi _+`$, this operation can be called ”abelian” twisting. Abelian because only phases are controlled. Thus we can summarize our considerations by the following statement. A two-qubit state has perfectly secure one bit of key with respect to basis $`\|ij`$, if and only if it is a twisted EPR state (by abelian twisting of Eq. (99)): $$\rho _{AB}=U\|\psi _+\psi _+\|U^{}.$$ (100) Moreover, if a state satisfies security condition approximately, it must be close in fidelity to some state $`U\psi _+`$. The quality of the bit of key is given by magnitude of a c-number $`\|\rho _{0011}\|`$. ### VI-B The general case. In this section we will represent in terms of Eve’s states the state which has both key part and shield. We will see then, how the twisting becomes ”nonabelian”, and the condition of closeness to pure state $`U\psi _+`$ changes into that of closeness to pbit. If we write state in basis of system $`AB`$ (key part) we get blocks $`A_{iji^{}j^{}}`$ instead of matrix elements $$\rho _{ABA^{}B^{}}=\underset{iji^{}j^{}}{}\|ij_{AB}i^{}j^{}\|A_{A^{}B^{}}^{iji^{}j^{}}.$$ (101) After suitable transformations (see Appendix XV-A for details), we arrive at the following form: $$\rho _{ABA^{}B^{}}=\underset{iji^{}j^{}}{}\sqrt{p_{ij}p_{i^{}j^{}}}\|ij_{AB}i^{}j^{}\|\left(𝒰_{ij}\sqrt{\rho _E^{ij}}\sqrt{\rho _E^{i^{}j^{}}}𝒰_{i^{}j^{}}^{}\right)^T,$$ (102) where the operator $`𝒰_{ij}^{}U_{ij}𝒲V_{ij}`$ maps the space $`_{A^{}B^{}}`$ exactly onto a support of $`\rho _E^{i^{}j^{}}`$ in space $`_E`$, and the dual operator $`𝒰_{ij}=V_{ij}^{}𝒲^{}U_{ij}^{}`$ maps the support of $`\rho _E^{ij}`$ back to $`_{A^{}B^{}}`$. Let us note, that in parallel to Eq. (93) we have that the trace norms of the blocks $`A`$ are connected with fidelities between Eve’s states $$A_{iji^{}j^{}}=\sqrt{p_{ij}p_{i^{}j^{}}}F(\rho _E^{ij},\rho _E^{i^{}j^{}})$$ (103) #### The case of two qubit key part If the key part is two qubit system we get $$\rho _{ABA^{}B^{}}=\left[\begin{array}{cccc}p_{00}[𝒰_{00}\rho _E^{00}𝒰_{00}^{}]^T& \times & \times & \sqrt{p_{00}p_{11}}[𝒰_{00}\sqrt{\rho _E^{00}}\sqrt{\rho _E^{11}}𝒰_{11}^{}]^T\\ \times & p_{01}[𝒰_{01}\rho _E^{01}𝒰_{01}^{}]^T& \times & \times \\ \times & \times & p_{10}[𝒰_{10}\rho _E^{10}𝒰_{10}^{}]^T& \times \\ \times & \times & \times & p_{00}[𝒰_{11}\rho _E^{11}𝒰_{11}^{}]^T\end{array}\right]$$ (104) Let us now discuss conditions for presence of one bit of key. They are again (i) $`p_{00}+p_{11}=\frac{1}{2}`$ and (ii) Eve’s states are the same $`\rho _E^{00}=\rho _E^{11}`$. This is equivalent to $$\sqrt{p_{00}p_{11}}F(\rho _E^{00},\rho _E^{11})=\frac{1}{2}$$ (105) which is nothing but trace norm of upper-right block $`A_{0011}`$. Also conditions for approximate bit of key requires the norm to be close to $`\frac{1}{2}`$. Moreover, to see how pbit and the twisting arise, let us put all Eve’s states equal to each other, and probabilities corresponding to perfect correlations. We then obtain $$\rho _{ABA^{}B^{}}=\frac{1}{d}\underset{ij}{}\|ii_{AB}jj\|[𝒰_{ii}\rho _E𝒰_{jj}^{}]^T.$$ (106) where $`\rho _E`$ is one fixed state, that Eve has irrespectively of outcomes. We see here almost the form of pbit. One difference might be is that instead of usual unitaries, we have some embeddings $`𝒰_{ii}`$. However, since now Eve’s space is of the same dimension as $`A^{}B^{}`$ (because Eve has single state), they are actually usual unitaries. The transposition does not really make a difference, as it can be absorbed both by state, and by unitaries. It is interesting to see here in place of phases from previous section the unitaries appeared, so that abelian twisting changed into nonabelian one. Also the condition for key changed from modulus of c-number - matrix element, to a trace norm of q-number - a block. ## VII Overview In this section we will shortly summarize what we have done so far. Then we will describe the goals of the paper, and briefly outline how we will achieve them. ### VII-A Pbits and twisting We have considered a state shared by Alice and Bob, which was divided into two parts: the key part $`AB`$ and the shield $`A^{}B^{}`$. The key part is measured in a local basis, while the shield is kept. The latter is seen by Eve as an environment that may restrict her knowledge about outcomes of measurement performed on the pdit. We have shown two important facts. First, we have characterized all the states for which measurement on the key part gives perfect key. The states are called pdits, and they have a very simple form. Moreover, we have shown that twisting does not change the ccq state arising from measurement on the key part part. (We should emphasize here, that twisting must be controlled by just the same basis in which the measurement is performed.) This is an interesting feature, because twisting may be a nonlocal transformation. Thus even though we apply a nonlocal transformation to the state, the quality of the key established by measuring the key part (in the same basis) does not change. ¿From the exhibited examples of pbits, we have seen that some of them have very small distillable entanglement. Since pbits are EPR states subjected to twisting, we see that in this case the twisting must have been very nonlocal, since it significantly diminished distillable entanglement. Because pbits contain at least one bit of secure key, we have already seen that distillable key can be much larger than distillable entanglement. However our main goal is to show that there are bound entangled states from which one can draw key. Thus we need distillable entanglement to be strictly zero. Here it is easily seen that any perfect pdit is an NPT state. Even more, one can show that pdits are always distillable . Thus we cannot realize our goals by analysing perfect pbits. ### VII-B Approximating pbits with PPT states After realising that pbits cannot be bound entangled, one finds that this still does not exclude bound entangled states with private key. Namely, even though bound entangled states cannot contain exact key (as they would be pbits then) they may contain almost exact key. Such states would be in some sense close to pbits. Note that this would be impossible, if the only states containing perfect key were maximally entangled state. Indeed, for $`dd`$ system if only a state has greater overlap than $`1/d`$ with a maximally entangled state we can distill singlets from it . Recall, that for a state with key part being two qubits, the measure of quality of the bit of key coming from measuring the key part is trace norm of upper-right block. The key is perfect if the norm is $`1/2`$ (we have then pbit) and it is close to perfect, if the norm is close to $`1/2`$. Thus our first goal will be to find bound entangled states having the trace norm of that block arbitrarily close to $`1/2`$. We will actually construct such PPT states (hence bound entangled) in Sections X-A, X-B. In this way we will show that there exist bound entangled states that contain an arbitrarily (though not perfectly) secure bit of key. ### VII-C Nonzero rate of key from bound entangled states It is not enough to construct bound entangled states with arbitrary secure single bit of key. The next important step is to show that given many copies of BE states one can draw nonzero asymptotic rate of secure key. To show this we will employ (in Section X-C) the BE states with almost perfect bit of key. Let us outline here the most direct way of proving the claim. To be more specific, we will consider many copies of states $`\rho _ϵ`$ which have upper-right block trace norm equal to $`1/2ϵ`$. We will argue that one can get key by measuring the key part of each of them, and then process via local classical manipulations and public discussion the outcomes. How to see that one can get nonzero rate in this way? We will first argue, that the situation is the same, as if the outcomes were obtained from a state which is close to maximally entangled. To this end we will apply the idea of privacy squeezing described in Section II-B. First recall, that we have shown that operation of twisting does not change security of ccq state - more precisely, it does not change the state of the Eve’s system and key part of Alice and Bob systems, which would arise, if Alice and Bob measured the key part. Thus whatever twisting we will apply, from cryptographic point of view the situation will not change. The total state will change, yet this can be noticed only by those who have access to the shield of Alice and Bob systems, and Eve does not have such access. We will choose such a twisting, that will change the upper-right block of $`\rho ^ϵ`$ into a positive operator. This is exactly the one which realizes privacy squeezing of this state. Now, even though security is not changed, the state is changed in a very favorable way for our purposes. Namely, we can now trace out the shield, and the remaining state of the key part (a p-squeezed state of the initial one) will be close to maximally entangled. Indeed, twisting does not change trace norm of the upper-right block. Because now the block is a positive operator, its trace norm is equal to its trace, and tracing out shield amounts just to evaluating trace of blocks. Since the trace norm was $`1/2ϵ`$, the upper-right element of the state of key part is is equal to $`1/2ϵ`$, which means that state is close to maximally entangled (where the corresponding element is equal to $`1/2`$). One can worry, that it is now not guaranteed that the security is the same, because we have performed not only twisting, but also partial trace over shield. However the latter operation could only make situation worse, since partial trace means giving the traced system to Eve. Now the only remaining thing is to show that we can draw key from data obtained by measuring many copies of state close to an EPR state, then definitely we can draw key from many copies of more secure ccq state obtained from our $`\rho ^ϵ`$. To achieve the goal, we thus need some results about drawing key from ccq state. Let us recall that we work in scenario, where Alice and Bob are promised to share i.i.d. state, so that the ccq state is tensor product of identical copies. In this case, the needed results have been provided in . It follows, that the rate of key is at least $`I(A:B)I(A:E)`$, where $`I`$ is mutual information. If instead of almost-EPR state, we have just an EPR state, the above quantity is equal to $`1`$. Indeed, perfect correlations, and perfect randomness of outcomes gives $`I(A:B)=1`$ and purity of the EPR state gives $`I(A:E)=0`$. Since we have state close to an EPR state, due to continuity of entropies, we will get $`I(A:B)1`$, and $`I(A:E)0`$. Thus given $`n`$ copies of states that approximate pbits one can get almost $`n`$ bits of key in limit of large $`n`$. ### VII-D Drawing key and transforming into pbits by LOCC Apart from showing that key can be drawn from BE states, we want to develop the theory of key distillation from quantum states. To this end in Section VIII we recast definition of distilling key in terms of distilling pbits by local operations and classical communication. This is important change of viewpoint: drawing key requires referring to Eve; while distilling pbits by LOCC concerns solely bipartite states shared by Alice and Bob, and never requires explicit referring to Eve’s system. Thus we are able to pass from the game involving three parties: Alice, Bob and Eve to the two players game, involving only Alice and Bob. We will employ two basic tools: (i) the concept of making a protocol coherent; (ii) the fact (which we will prove) that having almost perfectly secure ccq state is equivalent to having a state close to some pdit. Note that in one direction, the reasoning is very simple: if we can get nonzero rate of asymptotically perfect pdits by LOCC, we can also measure them at the end, and get in this way asymptotically perfect ccq states, which is ensured by item (ii) above. The converse direction is a little bit more involved: we take any protocol that produces key, apply it coherently, and this gives pure final state of Alice, Bob and Eve’s systems. From (ii) it follows, that the total state of Alice and Bob must be close to pdit. Let us briefly discuss how we will show the fact (ii). The essential observation is that both the ccq state $`\stackrel{~}{\rho }_{ABE}`$ and Alice and Bob total state $`\rho _{ABA^{}B^{}}`$ are reductions of the same pure state $`\psi _{ABA^{}B^{}E}`$. Here some explanation is needed: in general, since the ccq state is obtained by measurement, it is not reduction of $`\psi _{ABA^{}B^{}E}`$. However, one can first apply measurement coherently to the state $`\rho _{ABA^{}B^{}}`$. Then the ccq of the new state $`\rho _{ABA^{}B^{}}^{}`$ is indeed the reduction of $`\psi _{ABA^{}B^{}E}^{}`$. In the actual proof we will proceed in a slightly different way. Now, if we have two nearby ccq states, we can find their purifications that are close to each other too. Then also Alice and Bob states arising when we trace out Eve’s system are close to each other (because partial trace can only make states closer). The whole argument is slightly more complicated, but the above reasoning is the main tool. The equivalence we obtain puts the task of drawing key into the standard picture of state manipulations by means of LOCC. The theory of such manipulations is well developed and, in particular, there are quite general methods of obtaining bounds on transition rates (in our case the transition rate is just distillable key), see . Indeed, we will be able to show that relative entropy of entanglement is an upper bound for distillable key. The main idea of deriving the bound is similar to the methods from LOCC state manipulations. However significant obstacles arise, to overcome which we have developed essentially new tools. ## VIII Two definitions of distillable key: LOCC and LOPC paradigms In this section we show that distillable amount of pdits by use of LOCC denoted by $`K_D`$ is equal to classical secure key distillable by means of local operations and public communication (LOPC). ### VIII-A Distillation of pdits We have established a family of states - pdits - which have the following property: after measurement in some basis $``$ they give a perfect dit of key. In entanglement theory one of the important aims is to distill singlets (maximally entangled states) which leads to operational measure of distillable entanglement. We will pose now an analogous task namely distilling pdits (private states) which are of the form (6). This gives rise to a definition of distillable key i.e. maximal achievable rate of distillation of pdits. Similarly as in the case of distillation of singlet, it is usually not possible to distill exact pdits. Therefore the formal definition of distillable key $`K_D`$ will be a bit more involved. ###### Definition 6 For any given state $`\rho _{AB}(_A_B)`$ let us consider sequence $`P_n`$ of $`LOCC`$ operations such that $`P_n(\rho _{AB}^n)=\sigma _n`$, where $`\sigma _nB(_A^{(n)}_B^{(n)})`$. A set of operations $`𝒫_{n=1}^{\mathrm{}}\{P_n\}`$ is called pdit distillation protocol of state $`\rho _{AB}`$ if there holds $$\underset{n\mathrm{}}{lim}\sigma _n\gamma _{d_n}=0,$$ (107) where $`\gamma _{d_n}`$ is a pdit whose key part is of dimension $`d_n\times d_n`$. For given protocol $`𝒫`$, its rate is given by $$(𝒫)=\underset{n\mathrm{}}{lim\; sup}\frac{\mathrm{log}d_n}{n}$$ (108) The distillable key of state $`\rho _{AB}`$ is given by $$K_D(\rho _{AB})=\underset{𝒫}{sup}(𝒫).$$ (109) In other words, due to this definition, Alice and Bob given $`n`$ copies of state $`\rho _{AB}`$ try to get a state which is close to some pdit state with $`d=d_n`$. Unlike so far in entanglement theory, effect of distillation of quantum key depends not only on the number $`n`$ of copies of initial state but also on the choice of the output state. This is because private dits appears not to be reversibly transformable with each other by means of LOCC operations, as it is in case of maximally entangled states in LOCC entanglement distillation. Thus the quantity $`K_D`$ is a rate of distillation to the large class of states. (Of course, since the definition involves optimization, $`K_D`$ is well defined; in particular the expensive pdits will be suppressed). One can be interested now if this new parameter of states $`K_D(\rho )`$ has an operational meaning for quantum cryptography. One connection is obvious: given a quantum state Alice and Bob may try to distill some pdit state, and hence get (according to the above definition) $`K_D(\rho _{AB})`$ bits of key if such distillation has nonzero rate. However the question arises: is it the best way of extraction of a classical secure key from a quantum state? I.e. given a quantum state is the largest amount of classical key distillable from a state equal to $`K_D`$. We will give to this question a positive answer now. It means, that distilling private dits i.e. states of the form (22) is the best way of distilling classical key from a quantum state. ### VIII-B Distillable classical secure key: LOPC paradigm The issue of drawing classical secure key from a quantum state is formally quite different from the definition of drawing pdits. However it will turn out that it is essentially the same thing. In the LOCC paradigm, we have an initial state $`\rho `$ hold by Alice and Bob who apply to it an LOCC map, and obtain a final state $`\rho ^{}`$. Thus the LOCC paradigm is essentially a bipartite paradigm. In the paradigm of drawing secure classical key (see e.g. ), there are three parties, Alice, Bob and Eve. They start with some joint state $`\rho _{ABE}`$ where subsystems $`A,B,E`$ belong to Alice, Bob and Eve respectively. Now, Alice and Bob essentially perform again some LOCC operations. However we have now a tripartite system, and we should know how that operation act on the whole system. ###### Definition 7 An operation $`\mathrm{\Lambda }`$ belongs to LOPC class if it is composition of * Local Alice (Bob) operations, i.e. operations of the form $`\mathrm{\Lambda }_AI_{BE}`$ (or $`\mathrm{\Lambda }_BI_{AE}`$). * Public communication from Alice to Bob (and from Bob to Alice). E.g. the process of communication from Alice to Bob is described by the following map $$\mathrm{\Lambda }(\rho _{aABE})=\underset{i}{}P_i\rho _{aABE}P_i\|i_bi\|\|i_ei\|$$ (110) where $`P_i=I_{ABE}\|i_ai\|`$. Here the subsystem $`a`$ carries the message to be sent, and the subsystems $`b`$ and $`e`$ of Bob and Eve represent the received message. Let us note, that LOCC operations can be defined in the same way, the only difference is that we drop Eve’s systems (both $`E`$ and $`e`$) and corresponding operators. Now, drawing secure key means obtaining the following state $$\rho _{ideal}^{ccq}=\underset{i=0}{\overset{d1}{}}\frac{1}{d}\|iiii\|_{AB}\rho ^E$$ (111) by means of LOPC. Since output states usually can not be exactly $`\rho _{ideal}^{ccq}`$ Alice and Bob will get state of the ccq form (2) i.e. $$\rho _{real}^{ccq}=\underset{i,j=1}{\overset{d}{}}p_{ij}\|ij_{AB}ij\|\rho _{ij}^E$$ (112) There are two issues here: first, Alice and Bob should have almost perfect correlations, second, Eve states should have small correlations with states $`\|ij`$ of Alice and Bob systems. The first condition refers to uniformity, the second one to security. There are several ways of quantifying these correlations, and some of them are equivalent. To quantify security one can use Holevo function of distilled ccq state, namely: $$\chi (\rho _{ccq})S(\rho _E)\underset{i,j=1}{\overset{d}{}}p_{ij}S(\rho _{ij})ϵ$$ (113) where $`S(\rho )=\mathrm{Tr}\rho \mathrm{log}\rho `$ denotes von Neumann entropy, and $$\rho _E=\underset{i,j=1}{\overset{d}{}}p_{ij}\rho _{ij}.$$ (114) Alternatively, one can use similar condition based on norm $$\underset{ij}{}p_{ij}\rho _E\rho _{ij}^Eϵ$$ (115) The condition of maximal correlations between Alice and Bob (uniformity) can be of the following form $$\underset{i,j=1}{\overset{d}{}}p_{ij}\|ijij\|\frac{1}{d}\underset{i=1}{\overset{d}{}}\|iiii\|ϵ$$ (116) One can also use again the trace norm between the real state (112) that is obtained and the ideal desired state (111) as done in , which includes both maximal correlations condition as well as security condition. The condition says that the state $`\rho _{real}^{ccq}`$ obtained by Alice and Bob is closed to some ideal state $$\rho _{real}^{ccq}\rho _{ideal}^{ccq}ϵ$$ (117) We will discuss relations between this condition, and security criteria (113) and (115) as well as with uniformity criterion (116) in Appendix XV-C . For the purpose of definition of secret key rate in this paper, we apply the joint criterion (117). Consequently, we adopt the following measure of distillable classical secure key from a quantum tripartite state: ###### Definition 8 For any given state $`\rho _{ABE}B(_A_B_C)`$ let us consider sequence $`P_n`$ of $`LOPC`$ protocols such that $`P_n(\rho _{ABE}^n)=\beta _n^{}`$, where $`\beta _n^{}`$ is ccq state $$\beta _n^{}=\underset{i,j=0}{\overset{d_n1}{}}p_{ij}\|ijij\|_{AB}\rho _{ij}^E$$ (118) from $`(^{(n)})=B(_A^{(n)}_B^{(n)}_E^{(n)})`$ with $`dim_A^{(n)}=dim_B^{(n)}=d_n`$. A set of operations $`𝒫_{n=1}^{\mathrm{}}\{P_n\}`$ is called classical key distillation protocol of state $`\rho _{AB}`$ if there holds $$\underset{n\mathrm{}}{lim}\beta _n^{}\beta _{d_n}=0,$$ (119) where $`\beta _{d_n}B(^{(n)}`$ is of the form $$\frac{1}{d_n}\left(\underset{i=1}{\overset{d_n}{}}\|ii_{AB}ii\|\right)\rho _n^E,$$ (120) $`\rho _n^E`$ are arbitrary states from $`B(_E^{(n)})`$. The rate of a protocol $`𝒫`$ is given by $$(𝒫)=\underset{n\mathrm{}}{lim\; sup}\frac{\mathrm{log}d_n}{n}$$ (121) Then the distillable classical key of state $`\rho _{ABE}`$ is defined as supremum of rates $$C_D(\rho _{ABE})=\underset{𝒫}{sup}(𝒫).$$ (122) The above definition works for any input tripartite state $`\rho _{ABE}`$. However in this paper we are only interested in the case where the total state is pure. The latter is determined by state $`\rho _{AB}=\mathrm{Tr}_E\rho _{ABE}`$ up to unitary transformations on Eve’s side. Since from the very definition $`C_D`$ does not change under such transformations, the latter freedom is not an issue, so that we can say the state $`\rho _{AB}`$ completely determines the total state. Thus we get definition of distillable classical secure key from bipartite state $`\rho _{AB}`$ ###### Definition 9 For given bipartite state $`\rho _{AB}`$ the distillable classical secure key is given by $$C_D(\rho _{AB})C_D(\psi _{ABE})$$ (123) where $`\psi _{ABE}`$ purification of $`\rho _{AB}`$. ### VIII-C Comparison of paradigms Let us compare two definitions 109 and 9 of distilling cryptographical key. The difference is mostly that the first one deals only with bipartite system, and the goal is to get the desired final state by applying a class of LOCC operations. Within the second paradigm, we have tripartite state and we want to get a wanted state by means of LOPC operations. Thus the first paradigm is much more standard in quantum information theory. The second one comes from classical security theory (see e.g. ), where probability distributions of triples of random variables $`P(X,Y,Z)`$ are being processed. In the next section we will see that if the tripartite initial state is pure, the two paradigms are tightly connected. In the case of distillation of exact key, they are almost obviously identical, while in the inexact case, the only issue is to make the asymptotic security requirements equivalent. We will see that an output pdit obtained by LOCC implies some ccq state obtained by LOPC, and vice versa. ### VIII-D Composability issues In the present paper we consider the promised scenario, which is the first step to consider in unconditionally secure QKD. <sup>1</sup><sup>1</sup>1The next step was subsequently done in papers . The latter security definition is required to be universally composable, which means that a QKD protocol can be used as a subroutine of any other cryptographic protocol . To this end, one needs to choose carefully the measure of security. Our starting point is the LOPC paradigm, where the measure of security is the trace norm. This is compatible with , where it was shown that such security measure implies indeed composability (see eq. (10) of ). In particular, if we concatenate $`n`$ QKD protocols with security $`ϵ`$ measured by trace norm between the obtained state and the ideal target, the overall protocol will have security bounded by $`nϵ`$. Thus trace norm in LOPC paradigm can be called composable security condition. However, we want to recast QKD within the LOCC paradigm, i.e. in terms of distance between Alice and Bob states rather than tripartite states. In this spirit, in it is shown that fidelity with maximally entangled state implies composability, in the sense that if fidelity is $`1ϵ`$ then the norm is bounded by $`\sqrt{ϵ}`$ (see eq. (20) of ). This may seem a bit uncomfortable, because while composing many protocols we have now to add not epsilons, but rather their square roots. However one can see that it is in general not possible to do better while starting from fidelity. Due to inequality $`1F(\sigma ,\rho )\frac{1}{2}\rho \sigma `$ we can equally well use the trace norm. Indeed it will give us at worse the $`\sqrt{ϵ}`$ bound for composable security condition in the LOPC paradigm. In our paper we have a more general situation, i.e. we deal with distance to private states rather than just maximally entangled states. According to the above discussion, we have decided to use trace norm in LOCC scenario. In our techniques, in the mid-steps we shall use fidelity, hence we are again left with $`\sqrt{ϵ}`$. We do not know, whether one can omit fidelity, and get rid of the square root. Having said all that, we should emphasize that finally the minimal requirement is that security is an exponentially decreasing function of some security parameter, whose role can play e.g. the number of all qubits in the game. This condition is of course not spoiled by square root. (Example of a situation, where this requirement is not met was conjectured in and then proved in ). In our paper we shall derive equivalences between various security conditions, and the word “equivalence” will be understood in the sense of not spoiling the proper exponential dependence on total number of qubits. In particular, the security conditions will be called equivalent also when they differ by the factor polynomial in number of qubits (see e.g. ). ## IX Equality of key rates in LOCC and LOPC paradigms In this section we will show that definitions 109 and 9 give rise to the same quantities. In this way the problem of drawing key within original LOPC paradigm is recast in terms of transition to a desired state by LOCC. First we will describe a coherent version of LOPC protocol. Then we will use it to derive equivalence in exact case (where protocols produce as outputs ideal ccq states or ideal pdits). Subsequently we will turn to the general case where inexact transitions are allowed. ### IX-A Coherent version of LOPC key distillation protocol The main difference between LOPC and LOCC paradigms is that in the first one we have transformations between tripartite states shared by Alice, Bob and Eve, while in the latter one - between bipartite states shared by Alice and Bob. Thus in LOPC paradigm, the part of the state held by Alice and Bob does not, in general tell us about security. To judge if Alice and Bob have secure key we need the whole $`\rho _{ABE}`$ state. Security is assured by the lack of correlations of this state with Eve. Thus if we want to recast the task of drawing key in terms of LOCC paradigm, we need to consider such LOCC protocol which produces output state that assure security of key itself. We will do this by considering coherent version of LOPC key distillation protocols (cf. ). The most important feature of the version will be that given any LOPC protocol, starting with some initial pure state $`\psi _{ABE}`$ and ending up with some ccq state $`\rho _{ABE}`$, its coherent version will end up with a state $`\psi _{AA^{}BB^{}E}^{}`$ such that tracing out $`A^{}B^{}`$ part will give exactly the ccq state $`\rho _{ABE}`$. In this way, the total Alice and Bob state $`\rho _{AA^{}BB^{}}`$ will keep the whole information about Eve (because up to unitary on Eve’s system, purification is unique). In coherent version of key distillation protocol Alice and Bob perform their local operations in a coherent way i.e. by adding ancillas, performing unitary transformations and putting aside appropriate parts of the system. This additional part of the system is discarded in the usual protocol. However, holding this part allows one to keep the total state of Alice Bob and Eve pure in each step of the protocol. This is because we use pure ancillas, pure initial state and apply only unitary transformations which preserve purity. Alice and Bob can also perform public communication. Its coherent version is that e.g. Bob and Eve apply C-NOT operations to Alice’s subsystem which holds the result of measurement. Formally, the coherent version of process of communication from definition 7 is an operation of the form $$\mathrm{\Lambda }(\rho _{aABE})=U(\rho _{aABE}\|0_a^{}0\|\|0_b0\|\|0_e0\|)U^{}$$ (124) where $$U=I_{ABE}\underset{i}{}\|i_ai\|U_a^{}^{(i)}U_b^{(i)}U_e^{(i)}$$ (125) with unitary transformation $`U_e^{(i)}`$,$`U_b^{(i)}`$ satisfying $`U_e^{(i)}\|0_e=\|i_e`$ and similarly for $`U_b^{(i)}`$, $`U_a^{}^{(i)}`$. To finalize the operation, Alice puts aside the system $`a^{}`$. Such a coherent version of LOPC protocol has the following two features: (i) keeps the state pure. (ii) after tracing out subsystems that are put aside we obtain exactly the same state as in the original protocol. Now we are in position to construct for a given LOPC protocol a suitable LOCC protocol, which will output pbits when the former protocol will output ideal key. Namely, the local operations from LOPC protocol, we replace with their coherent versions. The public communication we take in original - incoherent - form (of course, without broadcast to Eve since we now deal with bipartite states). One notes that such operation produces the same state of Alice and Bob as the one produced by coherent version of the original LOPC protocol traced out over Eve’s system. This is because, if we trace out Eve, then there is no difference between coherent and incoherent versions of public communication. Thus from an LOPC protocol we have obtained some LOCC protocol - a special one, where local systems are not traced out. In this way one gets a bridge which joins the two approaches, and shows that different definitions of distillable key are equivalent. In particular, suppose that the LOPC protocol produced the ideal ccq state. Then the output of LOCC protocol obtained as a coherent version of this protocol will produce a state which (due to theorem 2) must be a pdit. ###### Remark 3 Of course, the notion of coherent version need not concern just some LOPC protocol. Also an LOCC operation, that contained measurements and partial traces can be made coherent, which in view of the above considerations means simply, that the systems are not traced out, but only ”kept aside” and measurements are replaced by appropriate local unitaries. Actually, if we include all pure ancillas that will be added in the course of realizing the LOCC operation, the coherent version of the operation is nothing but a closed LOCC operation introduced for sake of counting local resources such as local information. One can think of the shield part, as being the state of all the lab equipment and quantum states, left over from the process of key distillation. ### IX-B Equivalence of paradigms: The case of exact key Here we will consider the ideal case, where the distillation of the key gives exactly the demanded output state. One can state it formally and observe: ###### Proposition 5 Let $`K_D^{exact}`$ and $`C_D^{exact}`$ denote optimal rates achievable by LOPC and LOCC protocols which as outputs have exact ccq states (111) and pdit states (22), respectively. Then for any state $`\rho _{AB}`$ we have $$K_D^{exact}(\rho _{AB})=C_D^{exact}(\rho _{AB})$$ (126) ###### Proof: If after LOPC protocol $`𝒫`$, Alice and Bob obtained exact $`d\times d`$ ccq state (111), then the coherent application of $`𝒫`$ due to theorem 2 and discussion of section IX-A will produce pdit of the same dimension. Conversely, if by LOCC Alice and Bob can get a pdit, then after measurement, again by theorem 2, they will obtain exact ccq state (111) of the same dimension. ### IX-C Distillation of classical key and distillation of pdits - equivalence in general (asymptotically exact) case We will prove here the theorem, which implies, that even in nonexact case, distillation of pdits from initial bipartite state by LOCC is equivalent to distillation of key by LOPC from initial pure state, that is purification of the bipartite state. This in turn means that the rates in both paradigms are equal. ###### Theorem 5 Let Alice and Bob share a state $`\rho `$ such that Eve has it’s purification. Then the following holds: if Alice and Bob can distill by LOPC operations a state such that with Eve’s subsystem it is ccq state i.e. of the form $$\rho _{ABE}=\underset{i,j=1}{\overset{d}{}}p_{ij}\|ijij\|_{AB}\rho _{ij}^E,$$ (127) with $`\rho _{ABE}^{ccq}\rho _{ideal}^{ccq}ϵ`$, then they can distill by LOCC operations a state $`\rho _{out}`$ which is close to some pdit state $`\gamma `$ in trace norm: $$\rho _{out}\gamma \sqrt{2ϵ},$$ (128) where the key part of a pdit $`\gamma `$ is of dimension $`d\times d`$. Conversely, if by LOCC they can get state $`\rho _{out}`$ satisfying $`\rho _{out}\gamma ϵ`$, then by LOPC they can get state $`\rho _{ccq}`$ satisfying $`\rho _{ABE}^{ccq}\rho _{ideal}^{ccq}\sqrt{2ϵ}`$ ###### Proof: The ”if” part of this theorem is proven as follows. By assumption Alice and Bob are able to get by some LOPC protocol $`𝒫`$ a ccq state $`\rho _{ABE}`$ satisfying $$\rho _{ABE}^{ccq}\rho _{ideal}^{ccq}ϵ.$$ (129) Now by equivalence between norm and fidelity (Eq. (11) of Appendix) we can rewrite this inequality as follows $$F(\rho _{ABE}^{ccq},\rho _{ideal}^{ccq})>1\frac{1}{2}ϵ.$$ (130) By definition of fidelity $$F(\rho ,\sigma )=\underset{\psi ,\varphi }{\mathrm{max}}\|\psi \|\varphi \|$$ (131) where maximum is taken over all purifications $`\psi `$ and $`\varphi `$ of $`\rho `$ and $`\sigma `$ respectively, we can fix one of these purification arbitrarily, and optimise over the other one. Let us then choose such a purification $`\psi _{ABA^{}B^{}E}`$ of $`\rho _{ABE}^{ccq}`$ which is the output of coherent application of the mentioned protocol $`𝒫`$. There exists purification $`\varphi _{ABA^{}B^{}E}`$ of $`\rho _{ideal}`$ such that it’s overlap with $`\psi `$ is greater than $`1\frac{1}{2}ϵ`$. Since the fidelity can only increase after partial trace applied to both the states, it will be still greater than $`1\frac{1}{2}ϵ`$ once we trace over Eve’s subsystem. Thus we have $$F(\rho _{ABA^{}B^{}}^\psi ,\sigma _{ABA^{}B^{}}^\varphi )>1\frac{1}{2}ϵ.$$ (132) where $`\sigma _{ABA^{}B^{}}^\varphi `$ and $`\rho _{ABA^{}B^{}}^\psi `$ are partial traces of $`\varphi `$ and $`\psi `$ respectively. The state $`\sigma _{ABA^{}B^{}}^\varphi `$ (partial trace of $`\varphi `$) comes from purification of an ideal state, and by the very definition it is some pdit state $`\gamma `$. At the same time, the state $`\rho _{ABA^{}B^{}}^\psi `$ (partial trace of $`\psi `$) is the one which is the output of coherent application of protocol $`𝒫`$. Thus by coherent version of $`𝒫`$ Alice and Bob can obtain state close to pdit which proves the ”if” part of the theorem. To obtain equivalence let us prove now the converse implication. The proof is a sort of ”symmetric reflection” of the proof of the previous part. This time we assume that there exists LOCC protocol starting with $`\rho `$, ending up with final state $`\rho _{out}`$ with key part od $`d\times d`$ dimension which is close to some pdit in norm i.e. $$\rho _{out}\gamma ϵ.$$ (133) Due to equivalence between fidelity and norm, we have $$F(\rho _{out},\gamma )1ϵ/2$$ (134) The total state after protocol is $`\psi _{ABA^{}B^{}E}`$, and if partially traced over Eve it returns $`\rho _{out}`$. Then we can find such $`\varphi `$, purification of $`\gamma `$, that $`F(\psi ,\varphi )>1ϵ`$. Now let Alice and Bob measure the key part and trace out the shield. Then out of $`\psi `$ we get some ccq state $`\rho _{out}^{ccq}`$. The same operation applied to $`\varphi `$ gives ideal ccq state (111) $`\rho _{ideal}^{ccq}`$. The operation can only increase the fidelity, so that $$F(\rho _{out}^{ccq},\rho _{ideal}^{ccq})1ϵ/2$$ (135) Returning to norms we get $$\rho _{out}^{ccq}\rho _{ideal}^{ccq}\sqrt{2ϵ}.$$ (136) ## X Distilling key from bound entangled states In this section we will provide a family of states. Then we will show that for certain regions of parameters they have positive partial transpose (which means that they are non-distillable). Subsequently, we shall show that out of the above PPT states one can produce, by an LOCC operation, states arbitrarily close to pbits (which also implies that they are entangled, hence bound entangled). More precisely, for any $`ϵ`$ we will find PPT states, from which by a LOCC protocol, one gets with some probability a state $`ϵ`$-close to some pbit. Since LOCC preserves the PPT property, this shows that pbits can be approximated with arbitrary accuracy by PPT states, in sharp contrast with maximally entangled states. We then show how to get from a state sufficiently close to a pbit with non vanishing asymptotic rate of key. We obtain it by reducing the problem to drawing key from states that are close to the maximally entangled state. ### X-A The new family of PPT states … Here we will present a family of states, and will determine the range of parameters for which the states are PPT. The idea of construction of the family is based on the so called hiding states found by Eggeling and Werner in . Let us briefly recall this result. In it was shown that one can hide one bit of information in two states by correlating the bit of information with a pair of states which are almost indistinguishable by use of LOCC operations, yet being almost distinguishable by global operations. The resulting state with the hidden bit is of the form: $$\rho _{hb}=\frac{1}{2}\|00\|_{AB}\rho _{hiding}^1+\frac{1}{2}\|11\|\rho _{hiding}^2$$ (137) In it was shown that there are separable states, which can serve as arbitrarily good hiding states. These states are $$\tau _1=(\frac{\rho _s+\rho _a}{2})^k,\tau _2=(\rho _s)^k,$$ (138) where $`\rho _s`$ and $`\rho _a`$ are symmetric and antisymmetric Werner states (56). The higher is the parameter $`k`$, the more indistinguishable by LOCC protocols the states become. We adopt the idea of hiding bits to hide entanglement. Namely instead of bits one can correlate two orthogonal maximally entangled states with these two hiding states and get the state: $$\rho _{he}=\frac{1}{2}\|\psi _+\psi _+\|_{AB}\tau _1^{A^{}B^{}}+\frac{1}{2}\|\psi _{}\psi _{}\|_{AB}\tau _2^{A^{}B^{}}$$ (139) Let us recall, that our purpose is to get the family of states which though entangled are not distillable, and can approximate pdit states. Then the choice of $`\rho _{he}`$ as a starting point has double advantage. First, because $`\tau _1`$ and $`\tau _2`$ are hiding, $`\rho _{he}`$ will not allow for distillation of entanglement by just distinguishing them. Second, the hiding states are separable, so they do not bring in any entanglement to the state $`\rho _{he}`$. However the state (139) is obviously NPT. Indeed, consider partial transposition of $`BB^{}`$ system. It is composition of partial transpositions of $`B`$ and $`B^{}`$ subsystems. If one applies it to the state (139), one gets $`\rho _{ABA^{}B^{}}^\mathrm{\Gamma }=(I_AT_BI_A^{}T_B^{})(\rho _{ABA^{}B^{}})=`$ $`=\left[\begin{array}{cccc}\frac{1}{2}(\frac{\tau _1+\tau _2}{2})^\mathrm{\Gamma }& 0& 0& 0\\ 0& 0& \frac{1}{2}(\frac{\tau _1\tau _2}{2})^\mathrm{\Gamma }& 0\\ 0& \frac{1}{2}(\frac{\tau _1\tau _2}{2})^\mathrm{\Gamma }& 0& 0\\ 0& 0& 0& \frac{1}{2}(\frac{\tau _1+\tau _2}{2})^\mathrm{\Gamma }\end{array}\right]`$ (144) where $`\mathrm{\Gamma }`$ denotes partial transposition over subsystem $`B^{}`$ (as partial transposition over B caused interchange of blocks of matrix of (139)). This matrix is obviously not positive for the lack of middle-diagonal blocks. To prevent this we admix a separable state $`\frac{1}{2}(\|0101\|+\|1010\|)\tau _2`$ with a probability $`(12p)`$, where $`p(0,\frac{1}{2}]`$. It’s matrix reads then $$\rho _{(p,d,k)}=\left[\begin{array}{cccc}p(\frac{\tau _1+\tau _2}{2})& 0& 0& p(\frac{\tau _1\tau _2}{2})\\ 0& (\frac{1}{2}p)\tau _2& 0& 0\\ 0& 0& (\frac{1}{2}p)\tau _2& 0\\ p(\frac{\tau _1+\tau _2}{2})& 0& 0& p(\frac{\tau _1+\tau _2}{2})\end{array}\right],$$ (146) In subscript we explicitly write the parameters on which this state depends implicitly: $`d=d_A^{}=d_B^{}`$ is the dimension of symmetric and antisymmetric Werner states used for hiding states (138) and $`k`$ is parameter of tensoring in their construction. We shall see, that for some range of $`p`$, almost every state of this family is a PPT state. We formalise it in the next lemma. ###### Lemma 5 Let $`\rho _aB(𝒞^d𝒞^d)`$ and $`\rho _bB(𝒞^d𝒞^d)`$ be symmetric and antisymmetric Werner states respectively, and let $`k`$ be such that $$\tau _1=(\frac{\rho _s+\rho _a}{2})^k,\tau _2=(\rho _s)^k$$ (147) holds. Then for any $`p(0,\frac{1}{3}]`$ and any $`k`$ there exists $`d`$ such that state (146) has positive partial transposition. More specifically, the state (146) is PPT if and only if the following conditions are fulfilled $`0<p{\displaystyle \frac{1}{3}}`$ (148) $`{\displaystyle \frac{1p}{p}}\left({\displaystyle \frac{d}{d1}}\right)^k`$ For the proof of this lemma see Appendix XV-B. ### X-B … can approximate pdits We have just established a family $`\rho _{(p,d,k)}`$ such that for certain $`p`$, $`k`$ and $`d`$ they are PPT states. We will then show, that by LOCC one can transform some of them to a state close to pbits. More precisely, for any fixed accuracy, we will always find $`p`$, $`d`$ and $`k`$ such that it is possible to reach pbit up to this accuracy, starting from some number of copies of $`\rho _{(p,d,k)}`$ and applying LOCC operations. Subsequently, we will show that one can always choose the initial states $`\rho _{(p,d,k)}`$ to be PPT. Since LOCC operations do not change PPT property, we will in this way show that there are PPT states that approximate pbits to arbitrarily high accuracy. We will first prove the following theorem. ###### Theorem 6 For any $`ϵ>0`$ and any $`p(\frac{1}{4},1]`$ there exist state $`\rho `$ from family of state $`\{\rho _{(p,d,k)}\}`$ (146) such that for some $`m`$ from $`\rho ^m`$ one can get by LOCC (with nonzero probability of success) a state $`\sigma `$ satisfying $`\sigma \gamma ϵ`$ for some private bit $`\gamma `$. ###### Proof: First of all let us notice that by theorem 4 it is enough to show, that one can transform $`\rho ^m`$ into a state $`\rho ^{}`$ which has sufficiently large norm of the upper-right block. Let Alice and Bob share $`m`$ copies of a state $`\rho `$ from the family (146). The number $`m`$ and parameters $`(p,k,d)`$ of this state will be fixed later. Now let Alice and Bob apply the well known recurrence protocol - ingredient of protocols of distillation of singlet states . Namely they take one system in state $`\rho `$ as source system, and iterate the following procedure. In $`i`$-th step they take one system in state $`\rho `$, and treat it as a target system. Let us remind that both systems have four subsystems $`A`$, $`B`$, $`A^{}`$ and $`B^{}`$. To distinguish the source and target system, the corresponding subsystems of a target system we call $`\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{A^{}},\stackrel{~}{B^{}}`$. On the source and target system they both perform a CNOT gate with a source at the $`A(B)`$ part of a source system and target at $`\stackrel{~}{A}(\stackrel{~}{B})`$ part of a target system for Alice (Bob) respectively. Then, they both measure the $`\stackrel{~}{A}`$ and $`\stackrel{~}{B}`$ subsystem of the target system in computational basis respectively, and compare the results. If the results agree, they proceed the protocol, getting rid of the $`\stackrel{~}{A}\stackrel{~}{B}`$ subsystem. If they do not agree, they abort the protocol. With nonzero probability of success they can perform this operation $`m1`$ times having each time the same source system, and some fresh target system in state $`\rho `$. That is they start with $`m`$ systems in state $`\rho `$ and in each step (upon success) they use up one system and pass to the next step. One can easily check, that the submatrices (blocks) of the state $`\rho _{(p,d,k)}^{rec}`$ which survives $`m1`$ steps of this recurrence protocol (which clearly happens with nonzero probability) are equal to the $`m`$-fold tensor power of the elements of initial matrix $`\rho _{(p,d,k)}`$: $$\rho _{(p,d,k)}^{rec}=\frac{1}{N}\left[\begin{array}{cccc}[p(\frac{\tau _1+\tau _2}{2})]^m& 0& 0& [p(\frac{\tau _1\tau _2}{2})]^m\\ 0& [(\frac{1}{2}p)\tau _2]^m& 0& 0\\ 0& 0& [(\frac{1}{2}p)\tau _2]^m& 0\\ [p(\frac{\tau _1\tau _2}{2})]^m& 0& 0& [p(\frac{\tau _1+\tau _2}{2})]^m\end{array}\right].$$ (149) where the normalisation is given by $$N=\mathrm{Tr}[\rho _{(p,d,k)}^{rec}]=2p^m+2(\frac{1}{2}p)^m.$$ (150) Let us consider the upper-right block $`\stackrel{~}{A}_{0011}`$ of the matrix (149) without normalisation. Norm of this block is equal to $`\stackrel{~}{A}_{0011}=\left({\displaystyle \frac{p}{2}}\right)^m({\displaystyle \frac{\rho _a\rho _s}{2}})^k\rho _{s}^{}{}_{}{}^{k}=`$ $`\left({\displaystyle \frac{p}{2}}\right)^m\left(2(12^k)\right)^m=p^m(12^k)^m.`$ (151) where second equality is consequence of the fact, that $`\rho _a`$ and $`\rho _s`$ have orthogonal supports which gives that $`\rho _s^k`$ is orthogonal to any term in expansion of $`(\frac{\rho _a\rho _s}{2})^k`$ but the one $`\frac{1}{2^k}\rho _s^k`$. Thus the result is equal to norm of $`[(\frac{\rho _a\rho _s}{2})^k\frac{1}{2^k}\rho _s^k]`$ (which is $`(1\frac{1}{2^k})`$) plus norm of the difference $`\|\frac{1}{2^k}\rho _s^k\rho _s^k\|`$ which gives the above formula. Thus the norm of the upper-right block $`A_{0011}`$ of the state (149) is given by $$A_{0011}=\frac{1}{N}\stackrel{~}{A}_{0011}=\frac{1}{2}(1\frac{1}{2^k})^m\frac{1}{1+(\frac{12p}{2p})^m}.$$ (152) We want now to see, if we can make the norm to be arbitrary close to $`1/2`$. (then by Lemma 4 the state will be arbitrary close to a pbit). Since $`p>\frac{1}{4}`$, we get that $`(\frac{12p}{2p})^m`$ converges to 0 with $`m`$. Although increasing $`m`$ diminishes the term $`(1\frac{1}{2^k})^m`$, we can first fix $`k`$ large enough, so that the whole expression (152) will be as close to $`\frac{1}{2}`$ as it is required. Now we have the following situation. We know that for $`p(\frac{1}{4},1]`$ if $`\frac{1}{2}(1\frac{1}{2^k})^m\frac{1}{1+(\frac{12p}{2p})^m}`$ is close to $`1/2`$, then the state (149) is close to pbit. On the other hand, from lemma (5) it follows that for (i) $`p(0,1/3]`$ and (ii) $`\frac{1p}{p}\left(\frac{d}{d1}\right)^k`$ the state $`\rho _{(p,d,k)}`$ is PPT, hence also the state (149) is PPT (because it was obtained from the former one by LOCC operation). If we now fix $`p`$ from interval $`(1/4,1/3]`$, then by choosing high $`m`$ and for such $`m`$, high enough $`k`$, then the state (149) is close to pbit. Now, we can fix also $`m`$ and $`k`$, and choose $`d`$ so large that the condition (ii) is also fulfilled so that the state becomes PPT. This proves the following theorem, which is main result of this section. ###### Theorem 7 PPT states can be arbitrarily close to pbit in trace norm. Here might be the appropriate place to note an amusing property of state (139): Namely, Eve knows one bit of information about Alice and Bob’s state – she knows the phase of their Bell state. But she only has one bit of information about their state, thus it cannot be that she also knows the bit of their state, which is the key. In some sense, giving Eve the bit of phase information, means that she cannot know the bit value. ### X-C Distillation of secure key In the previous subsection we have shown that private bits can be approximated by PPT states. Now, the question is whether given many copies of one of such PPT states Alice and Bob can get nonzero rate of classical key. Below, we will give the positive answer. The main idea of the proof is to show that from the PPT state which is close to pbit Alice and Bob by measuring, can obtain ccq state satisfying conditions of protocol (DW) found by Devetak and Winter . Namely, they have shown that for an initial cqq state (state which is classical only on Alice side; this includes ccq state as special case) between Alice, Bob and Eve, $$C_D(\rho _{ABE})I(A:B)I(A:E)$$ (153) Here $`I(A:B)`$ stands for the quantum mutual information of the state $`\rho `$ with subsystems $`A`$ and $`B`$ given by $$I(A:B)=S(A)+S(B)S(AB),$$ (154) where $`S(X)`$ stands for the von Neumann entropy of $`X`$ (sub)system of the state $`\rho `$. Using the above result we can prove now, that from many copies of states close to a pbit, one can draw nonzero asymptotic rate of key. ###### Lemma 6 If a state $`\rho `$ is close enough to pbit in trace norm, then $`K_D(\rho )>0`$. ###### Proof: The idea of the proof is as follows. Suppose that $`\sigma `$ is close to pbit $`\gamma `$. We then consider twisting that changes $`\gamma `$ into basic pbit $`P_{+}^{}{}_{AB}{}^{}\sigma _{A^{}B^{}}`$ with some state $`\sigma `$ on $`A^{}B^{}`$. We apply twisting to both states, so that they are still close to each other. Of course, this is only a mathematical tool: Alice and Bob cannot apply twisting, which is usually a nonlocal operation. The main point is that after twisting, according to theorem 1 the ccq state does not change. If we now trace out systems $`A^{}B^{}`$ the resulting state will be close to maximally entangled, and the resulting ccq state – at most worse from Alice and Bob point of view (because tracing out means giving to Eve). What we have done is just privacy squeezing of the state $`\rho `$. Now, the latter ccq state has come from measurement of a state close to the maximally entangled one. Thus the task reduces to estimate quantities $`I(A:B)`$ and $`I(A:E)`$ for a ccq state obtained from measuring the maximally entangled state. However due to suitable continuities, first one is close to $`1`$ and second one close to $`0`$. Now by DW protocol, one can draw a pretty high rate of key from such ccq state. Let us now proceed with the formal proof. We assume that for some pbit $`\gamma `$ we have $$\rho \gamma ϵ.$$ (155) Let us consider twisting $`U`$ which changes pdit $`\gamma `$ into a basic pdit. Existence of such $`U`$ is assured by theorem 2. If both states $`\gamma `$ and $`\rho `$ are subjected to this transformation, the norm is preserved, so that $$U\rho U^{}U\gamma U^{}ϵ.$$ (156) Also due to theorem 1 the ccq state obtained by measuring key part of $`U\rho U^{}`$ is the same as that from $`\rho `$. Now, the amount of key drawn from such ccq state will not increase if we trace out shield. Thus we apply such partial trace to $`U\rho U^{}`$ and to $`\gamma `$, and by monotonicity of trace norm get $$\stackrel{~}{\rho }_{AB}P_{AB}^+ϵ.$$ (157) It is now enough to show that from ccq state obtained by measuring $`\stackrel{~}{\rho }_{AB}`$ (where Eve holds the rest of its purification) one can get nonzero rate of key. To this end let us note, that for ccq state obtained from any bipartite state $`\rho _{AB}`$ by measuring its purification on subsystems A and B in computational basis, we have the following bound for $`I(A:E)`$: $$I(A:E)_{\mathrm{ccq}}S(\rho _{AB}).$$ (158) Now, since our state $`\stackrel{~}{\rho }_{AB}`$ is close to $`P_+`$, for which $`S(P_+)=0`$ and $`I(A:B)=1`$, we can use continuity of entropy, to bound these quantities for the state. ¿From Fannes inequality (see eq. 12 in Sec. II-A), we get $`I(A:B)_{\mathrm{ccq}}14ϵ\mathrm{log}d_{AB}h(ϵ),`$ (159) $`I(A:E)_{\mathrm{ccq}}S(\stackrel{~}{\rho }_{AB})8ϵh(ϵ).`$ (160) To get first estimate, it is enough to note, that due to monotonicity of trace norm, the estimate (157) is also valid if we dephase the state $`\stackrel{~}{\rho }_{AB}`$ and $`P_{AB}^+`$. Thus we obtain that $$K_D(\rho )I(A:B)I(A:E)116ϵ$$ (161) This ends the proof of the lemma. Let us note here, that it was not necessary to know, that the state $`\rho `$ is close to pbit. Rather, It was enough to know that trace norm of upper right block is close to $`1/2`$, as we have proved that it is equivalent to previous condition (see Sec. V). ¿From this it follows, that after twisting, and tracing out $`A^{}B^{}`$ the resulting state $`\stackrel{~}{\rho }`$ is close to the EPR state, which ensures nonzero rate of key (actually the rate is close to $`1`$). We now can combine the lemma with the fact that we know PPT states that are close to $`pbit`$, to obtain that there exist PPT states from which one can draw secure key. The states must be entangled, as from separable states one cannot draw key. Namely separable state can be established by public discussion. If it could then serve as a source of secret key, one could obtain secret key by public discussion which can not be possible. For formal arguments see . Thus our PPT states are entangled. But, since they are PPT, one cannot distill singlets from them , hence they are bound entangled. In this way we have obtained the following theorem ###### Theorem 8 There exist bound entangled states with $`K_D>0`$. We have split the way towards bound entangled states with nonzero key into two parts. First, we have shown that from PPT states $`\rho _{(p,d,k)}`$ by recurrence one can get a state that is close to pbit. Then we have shown, that from a state close to pbit one can draw private key. Note that we have two quite different steps: recurrence was the quantum operation preformed on quantum Alice and Bob states, while Devetak-Winter protocol in our case, is classical processing of the outputs of measurement. We could unify the picture in two ways. First Alice and Bob could measure the key part of the initial state $`\rho _{(p,d,k)}`$, and preform recurrence classically (since the quantum recurrence is merely coherent application of classical protocol). Then the whole process of drawing key from $`\rho _{(p,d,k)}`$ would be classical (of course, taking into account that Eve has quantum states). On the other hand, the DW protocol could be applied coheretly, so that till the very end, we would have quantum state of Alice and Bob. The result we have obtained allows to distinguish two measures of entanglement ###### Corrolary 2 Distillable entanglement and distillable classical secure key are different measures of entanglement i.e. there are states for which there holds $$K_D(\rho )>D(\rho )=0.$$ (162) In further section we will also show that $`K_D`$ is different than entanglement cost, as it is bounded by relative entropy of entanglement. ## XI Relative entropy of entanglement as upper bound on distillable key In this section we will provide complete proof of the theorem announced which gives general upper bound on distillable key $`K_D`$. This upper bound is given by regularised relative entropy of entanglement (62). The relative entropy of entanglement is given by $$E_r(\rho )=\underset{\sigma _{sep}}{inf}S(\rho \|\sigma _{sep}),$$ (163) where $`S(\rho \|\sigma )=\mathrm{Tr}\rho \mathrm{log}\rho \mathrm{Tr}\rho \mathrm{log}\sigma `$ is relative entropy, and infimum is taken over all separable states $`\sigma _{sep}`$. The regularized version of $`E_r`$ is given by $$E_r^{\mathrm{}}(\rho )=\underset{n}{lim}\frac{E_r(\rho ^n)}{n}.$$ (164) The limit exists, and due to subadditivity of $`E_r`$, we have $$E_r^{\mathrm{}}(\rho )E_r.$$ (165) It follows that also relative entropy of entanglement is upper bound for $`K_D`$. We recall now the following lemma obtained in , the proof of which we provide in Appendix XV-H: ###### Lemma 7 Consider a set $`𝒮^\tau :=\{U\rho _{ABA^{}B^{}}U^{}\|\rho _{ABA^{}B^{}}SEP,\rho _{ABA^{}B^{}}B(𝒞^d𝒞^d𝒞^{d_A^{}}𝒞^{d_B^{}})\}`$ where $`U`$ is $``$-twisting with $``$ being a standard product basis in $`𝒞^d𝒞^d`$. Let $`\sigma _{ABA^{}B^{}}𝒮^\tau `$ and $`\sigma _{AB}=\mathrm{Tr}_{A^{}B^{}}\sigma _{ABA^{}B^{}}`$. We have then $$S(P_+\|\sigma _{AB})\mathrm{log}d,$$ (166) where $`P_+=\|\psi _d^+\psi _d^+\|`$. We will also need asymptotic continuity of the relative entropy distance from some set of states obtained in in the form of . ###### Proposition 6 For any compact, convex set of state $`𝒮`$ that contains maximally mixed state, the relative entropy distance from this set given by $$E_r^𝒮=\underset{\sigma 𝒮}{inf}S(\rho \|\|\sigma ),$$ (167) is asymptotically continuous i.e. it satisfies $$\|E_r^𝒮(\rho _1)E_r^𝒮(\rho _2)\|<4ϵ\mathrm{log}d+h(ϵ)$$ (168) for any states $`\rho _1`$, $`\rho _2`$ acting on Hilbert space $``$ of dimension $`d`$, with $`ϵ=\rho _1\rho _2`$ with $`ϵ1`$. Let us mention, that the original relative entropy of entanglement has in place of $`𝒮`$ the set of separable states. Another version has been considered in , where $`𝒮`$ was set of PPT states. The latter set has entangled states, but they can be only weakly entangled. In contrast we will have set in which there may be quite strongly entangled states. We are now in position to formulate and prove the main result of this section. ###### Theorem 9 For any bipartite state $`\rho _{AB}B(𝒞^{d_A}𝒞^{d_B})`$ there holds $$K_D(\rho _{AB})E_r^{\mathrm{}}(\rho _{AB}),$$ (169) ###### Proof: By definition of $`K_D(\rho _{AB})`$ there exists protocol (i.e. sequence of maps $`\mathrm{\Lambda }_n`$), such that $$\mathrm{\Lambda }_n(\rho ^n)=\gamma _d^{}$$ (170) where $$\underset{n}{lim}\frac{\mathrm{log}d}{n}=K_D(\rho _{AB})$$ (171) and $$\underset{n}{lim}\gamma _d^{}\gamma _d\underset{n}{lim}ϵ_n=0$$ (172) with $`\gamma _d`$ being pdit with dimension $`d^2`$ of the key part. We will present now the chain of (in)equalities, and comment it below. $`S(\rho _{AB}^n\|\stackrel{~}{\sigma }_{sep})S(\gamma _{ABA^{}B^{}}^{}\|\sigma _{sep})=`$ (173) $`=`$ $`S(U_\gamma \gamma _{ABA^{}B^{}}^{}U_\gamma ^{}\|U_\gamma \sigma _{sep}U_\gamma ^{})`$ (174) $``$ $`S(\mathrm{Tr}_{A^{}B^{}}[U_\gamma \gamma _{ABA^{}B^{}}^{}U_\gamma ^{}]\|\mathrm{Tr}_{A^{}B^{}}[U_\gamma \sigma _{sep}U_\gamma ^{}])`$ (175) $``$ $`S(P_+^{}\|\sigma )`$ (176) $``$ $`inf_{\sigma T}S(P_+^{}\|\sigma ):=E_r^T(P_+^{})`$ (177) $``$ $`E_r^T(P_+)4P_+P_+^{}\mathrm{log}dh(P_+P_+^{})`$ (178) $``$ $`(14ϵ_n)\mathrm{log}dh(ϵ_n)`$ (179) Inequality (173) is due to the fact, that relative entropy does not increase under completely positive maps; in particular it can not increase under LOCC action applied to it’s both arguments (second argument becomes other separable state since LOCC operations can not create entanglement). In the next step, Eq. (174) we perform twisting $`U_\gamma `$ controlled by the basis in which state $`\gamma _m`$ is secure (without loss of generality we can assume it is standard basis). The equality follows from the fact that unitary transformation doesn’t change the relative entropy. Next (175) we trace out $`A^{}B^{}`$ subsystem of both states which only decreases the relative entropy. After this operation, the first argument is $`P_+^{}`$, which is a state close to the EPR state $`P_+`$. ($`P_+^{}`$ would be equal to the EPR state if $`\gamma _{ABA^{}B^{}}^{}`$ were exactly pdit) while second argument becomes some – not necessarily separable – state $`\sigma `$. The state belongs to the set $`T`$ constructed as follows. We take set of separable states on system $`ABA^{}B^{}`$ subject to twisting $`U_\gamma `$ and subsequently trace out the $`A^{}B^{}`$ subsystem. The inequality (176) holds, because we take infimum over all states from set $`T`$ of the function $`S(P_+^{}\|\sigma )`$. This minimised version is named there $`E_r^T(P_+^{})`$ as it is relative entropy distance of $`P_+^{}`$ from the set $`T`$. Let us check now, that set $`T`$ fulfills the conditions of proposition 6. Convexity of this set is obvious, since (for fixed unitary $`U_\gamma `$) by linearity it is due to convexity of the set of separable states. This set contains the identity state, since it contains maximally mixed separable state which is unitarily invariant (i.e. invariant under $`U_\gamma `$ ) and whose subsystem $`AB`$ by definition is the maximally mixed state as well. Thus by proposition 6 we have that $`E_r^T`$ is asymptotically continuous $$\|E_r^T(P_+^{})E_r^T(P_d^+)\|<P_+^{}P_+4\mathrm{log}d+h(P_+^{}P_+),$$ (180) where we assume that the EPR state $`P_+`$ is of local dimension $`d`$. Since $`P_+^{}`$ and $`P_+`$ come out of $`\gamma _{AB}^{}`$ and $`\gamma _d`$ by the same transformation described above (twisting, and partial trace) which doesn’t increase norm distance, by (172) we have that $`P_+^{}P_+ϵ_n`$. This, together with asymptotic continuity (180) implies (177). Now by lemma 7 we have $$E_r^T(P_+)\mathrm{log}d,$$ (181) which gives the last inequality: $$E_r^T(P_+^{})(14ϵ_n)\mathrm{log}dh(ϵ_n).$$ (182) Summarizing this chain of inequalities (173)-(178), we have that for any separable state $`\stackrel{~}{\sigma }_{sep}`$: $$S(\rho _{AB}^n\|\|\stackrel{~}{\sigma }_{sep})(14ϵ_n)\mathrm{log}dh(ϵ_n)$$ (183) Taking now infimum over all separable states $`\stackrel{~}{\sigma }_{sep}`$ we get $$E_r(\rho _{AB}^n)(14ϵ_n)\mathrm{log}dh(ϵ_n).$$ (184) Now we divide both sides by $`n`$ and take the limit. Then the left-hand-side converges to $`E_r^{\mathrm{}}`$. Due to (172) $`ϵ_n0`$ and due to (171), $`\mathrm{log}d/nK_D(\rho _{AB})`$. Thus due to continuity of $`h`$ we obtain $$E_r^{\mathrm{}}K_D$$ (185) As an application of the above upper bound, we consider now the relation between distillable key and entanglement cost. For maximally entangled states these two quantities are of course equal, unlike for general pdits. As an example let us consider again a flower state given in eq. (57). As follows from , the flower state has $`E_C`$ strictly greater than the relative entropy of entanglement. Since by the above theorem we have $`K_D(\gamma _{flower})E_r(\gamma _{flower})`$, havig $`E_r(\gamma _{flower})<E_C(\gamma _{flower})`$, we obtain in this case $`K_D(\gamma _{flower})<E_C(\gamma _{flower})`$. Let us note, that the entanglement monotone approach initiated here was then used in full extent in . It is shown there, that in fact any bipartite monotone $`E`$, which is continuous and normalized on private states (i.e. $`E(\gamma _d)\mathrm{log}d`$), is an upper bound on distillable key. In particular it is shown that the squashed entanglement is also an upper bound on distillable key. ## XII A candidate for NPT bound entanglement Thus far, all known bound entangled states have positive partial transpose (are PPT). A long-standing and interesting open question is whether there exist bound entangled states which are also NPT. If such states existed, it would imply that the quantum channel capacity is non-additive. Since any NPT state is distillable with the aid of some PPT state , we would have the curious property that one can have two states which are each non-distillable, but if you have both states, then the joint state would be distillable. We now present a candidate for NPT bound entangled states which are based on the states of equation (139), $$\rho _{he}=\frac{1}{2}\|\psi _+\psi _+\|_{AB}\tau _1^{A^{}B^{}}+\frac{1}{2}\|\psi _{}\psi _{}\|_{AB}\tau _2^{A^{}B^{}}$$ (186) and which intuitively appear to be bound entangled. Globally, the flags $`\tau _i^{A^{}B^{}}`$, are distinguishable, but under LOCC the flags appear almost identical, thus after Alice and Bob attempt to distinguish the flags, the state on $`AB`$ will be very close to an equal mixture of $`\psi _+`$ and $`\psi _{}`$. The equal mixture of only two different EPR states is separable in dimension $`2\times 2`$, but it is at the edge of separability. A slight biasing of the mixture, causes the state to be entangled. Thus, if Alice and Bob are able to obtain even a small amount of information about which $`\tau _i`$ they have, they will have a distillable state. More explicitly, if Alice and Bob attempt distillation by first guessing which hiding state flag they have, and then grouping the remaining parts of the states into two sets depending on their guess of the hiding state, they will be left with states of the form $$\rho _{he}=(\frac{1}{2}+ϵ)\|\psi _+\psi _+\|_{AB}+(\frac{1}{2}ϵ)\|\psi _{}\psi _{}\|_{AB}.$$ (187) This state is distillable. But what if we mix in more than two different EPR states? Namely, instead of only considering hiding states (flags) correlated to odd parity Bell states (anti-key states) $`\|\psi _\pm =\frac{1}{\sqrt{2}}(\|01\pm \|10)`$, we also add mix in flags correlated to the even parity Bell states (key type states) $`\|\varphi _\pm =\frac{1}{\sqrt{2}}(\|00\pm \|11)`$. Consider: $`\rho =p_{11}\|\varphi _+\varphi _+\|\rho _{11}+p_{12}\|\varphi _{}\varphi _{}\|\rho _{12}+`$ $`+p_{21}\|\psi _+\psi _+\|\rho _{21}+p_{22}\|\psi _{}\psi _{}\|\rho _{22}`$ (188) where $$\rho _{ij}=\tau _i\tau _j.$$ (189) Let us take for example, all $`p_{ij}=1/4`$. Then, after attempting to distinguish the hiding states, Alice and Bob will have a state which is very close to the maximally mixed state (i.e. the state will be very close to a mixture of all four Bell states). The maximally mixed state is very far from being entangled, thus even if Alice and Bob’s measurements on the hiding states are able to bias the mixture away from the maximally mixed state, the state will still be separable. Intuitively, it is thus clear why the state of equation (188) will not be distillable. Any protocol which attempts to first distinguish which Bell state the parties have, will fail. But is the state entangled? Indeed it is, in fact it has negative partial transpose. To see this, we look at the block-matrix form of the state $$\rho =\frac{1}{4}\left[\begin{array}{cccc}\tau _1(\tau _1+\tau _2)& 0& 0& \tau _1(\tau _1\tau _2)\\ 0& \tau _2(\tau _1+\tau _2)& \tau _2(\tau _1\tau _2)& 0\\ 0& \tau _2(\tau _1\tau _2)& \tau _2(\tau _1+\tau _2)& 0\\ \tau _1(\tau _1\tau _2)& 0& 0& \tau _1(\tau _1+\tau _2)\end{array}\right]$$ (190) If the matrix were PPT we would have in particular $$\tau _1^\mathrm{\Gamma }(\tau _1^\mathrm{\Gamma }+\tau _2^\mathrm{\Gamma })\tau _2^\mathrm{\Gamma }(\tau _1^\mathrm{\Gamma }\tau _2^\mathrm{\Gamma })$$ (191) We will argue that it is not true. Let us recall that $`\tau _1^\mathrm{\Gamma }=\left({\displaystyle \frac{P_+^{}}{d^21}}\right)^k`$ $`\tau _2^\mathrm{\Gamma }=\left({\displaystyle \frac{P_+^{}}{d^2+d}}+{\displaystyle \frac{(1+d)P_+}{d^2+d}}\right)^k\left({\displaystyle \frac{P_+^{}}{d^2+d}}\right)^k+R`$ $`\tau _1^\mathrm{\Gamma }\tau _2^\mathrm{\Gamma }=\left[(P_+^{})^k\left({\displaystyle \frac{1}{(d^21)^k}}{\displaystyle \frac{1}{(d^2+d)^k}}\right)R\right]`$ $`\tau _1^\mathrm{\Gamma }+\tau _2^\mathrm{\Gamma }=\left[(P_+^{})^k\left({\displaystyle \frac{1}{(d^21)^k}}+{\displaystyle \frac{1}{(d^2+d)^k}}\right)+R\right]`$ where we use notation from sec. XI. To see that (191) is not satisfied we consider the following projector $$Q=I(P_+^{})^k$$ (193) i.e. $`Q`$ is projector onto support of positive operator $`R`$. Since $`\mathrm{Tr}(Q\tau _1^\mathrm{\Gamma })=0`$, we have $$\mathrm{Tr}[(Q(P_+^{})^k)(\tau _1^\mathrm{\Gamma }(\tau _1^\mathrm{\Gamma }+\tau _2^\mathrm{\Gamma }))]=0$$ (194) Moreover we have $`\mathrm{Tr}(Q(P_+^{})^k)[\tau _2^\mathrm{\Gamma }(\tau _1^\mathrm{\Gamma }\tau _2^\mathrm{\Gamma })]=`$ $`=\mathrm{Tr}R\left({\displaystyle \frac{1}{(d^21)^k}}{\displaystyle \frac{1}{(d^2+d)^k}}\right)`$ (195) The above quantity is strictly greater than zero for $`d2`$. Thus inequality (191) is violated on projector $`Q(P_+^{})^k`$. Now, it may be that there is a protocol which succeeds in distilling from the state (188) which does not rely on first performing a measurement to distinguish the hiding states. However, even taking many copies of the state, produces a state of the form $$\rho =\underset{i}{}\|\psi _i\psi _i\|\rho _i$$ (196) with the $`\rho _i`$ being binary strings encoded in hiding states and $`\psi _i`$ being the basis of maximally entangled states. Thus the form of the state is invariant under tensoring. There is thus a very strong intuition that these states are NPT bound entangled, and a very good understanding of why they might be so. Effectively, the partial transpose does not feel very strongly the fact that the states $`\rho `$ are hiding states, but more strongly feels the fact that they are globally orthogonal. ## XIII Controlled private quantum channels Here, we demonstrate a cryptographic application of bound entangled states which have key. A private quantum channel (PQC) allows for the sending of quantum states such that an eavesdropper learns nothing about the sent states. Here, we consider the cryptographic primitive of having the ability to securely send quantum states (a PQC), but that this ability can be turned on and off by a controller. Namely, we consider a three party scenario (Alice, Bob, and the (C)controller) and demand * Alice and Bob have a private quantum channel, which they can use to send an unknown qubit from one to the other in such a way that they can be sure that no eavesdropper (including the Controller), can gain information about the state being sent. * the Controller has the ability to determine whether or not Alice and Bob can send the qubit We now show that this can be done using shared quantum states in such a way that the Controller only needs to send classical communication to one of the parties in order to activate the channel. First, let us note that the standard way of controlling the entanglement of two parties is via the GHZ state $$\|\psi _{ABC}=\|000+\|111.$$ (197) If the Controller, (Claire), measures in the basis $`\|0\pm \|1`$, then, depending on the outcome, Alice and Bob will share either the Bell state $`\|\psi _+`$ or $`\|\psi _{}`$. If $`C`$ then tells them the result, they will have one unit of entanglement (ebit) which they can then use to teleport quantum states. However, if the Controller wants to give them the ability to send a single qubit securely, then the GHZ state cannot be used for this, because the Controller can trick Alice and Bob into sending part of the quantum state to her. She can claim that she obtained measurement outcome $`+`$, when in reality she has not performed a measurement at all. Then, when Alice attempts to teleport a qubit to Bob, she is in fact teleporting to both Bob and the Controller. The controller can then perform a measurement on her qubit to obtain partial information about the sent qubit. Note that here we are concerned with the ability to give single shot access to a quantum channel. If the controller gives Alice and Bob many ebits by performing measurements on many copies of a GHZ state, then Alice and Bob could always perform purity testing to determine that the Controller is honest. Let us know show that unlike the GHZ, the states of Eq. (139) can be used in such a way that the Controller can give Alice and Bob single shot access to a private quantum channel, in such a way that Alice and Bob are sure that the Controller cannot obtain any information about the sent states even when the Controller cheats. We will then show that we can do the same thing with fully bound entangled states, so that Alice and Bob possess no distillable entanglement unless the Controller gives it to them. First, we assume the shared state as a trusted resource. I.e. a trusted party gives Alice, Bob and the Controller some state which they use to implement the primitive. This assumption can be removed in the limit of many copies, since if Alice and Bob have many copies of the state, they can perform tomography to ensure that they indeed possess the correct state. The state we initially use is the purification of Eq. (139) $$\rho _{he}=\frac{1}{2}\|\psi _+\psi _+\|_{AB}\tau _1^{A^{}B^{}}+\frac{1}{2}\|\psi _{}\psi _{}\|_{AB}\tau _2^{A^{}B^{}}$$ (198) Namely, $$\|\psi _{ABC}=\|00_{AB}\|\varphi _1_{A^{}B^{}C}+\|11_{AB}\|\varphi _2_{A^{}B^{}C}$$ (199) such that $`\mathrm{Tr}_C(\|\varphi _i\varphi _i\|)=\tau _i`$. Thus, $`\varphi _1\|\varphi _2=0`$ and since the $`\tau _i`$ are orthogonal, the Controller’s states $`\mathrm{Tr}_{A^{}B^{}}(\|\varphi _i\varphi _i\|)=\sigma _C^i`$ will be orthogonal. The controller can thus give Alice and Bob one ebit by performing a measurement to distinguish the $`\sigma _C^i`$. She then tells Alice and Bob the result. Alice and Bob on the other hand, are guaranteed security by the fact that they either possess the state $`\|\psi _+`$ or $`\|\psi _{}`$. I.e. it is an incoherent mixture of the two states, and they either have one of the states or the other, they just don’t know which one they have. The state of Eq. (139) however, does have an arbitrarily small amount of distillable entanglement. Thus, Alice and Bob will have access to a private quantum channel in the case of having many copies of the state. If we want to give full control to Claire, we need to ensure that the state held by Alice and Bob in the absence of Claire’s communication is non-distillable. This can be achieved by using the bound entangled states of equation (149) which approximate a pbit. It is not hard to verify, by explicitly writing the state in the Bell basis on $`AB`$, that the state is arbitrarily close to a state of the same form as equation (199), and thus has the desired properties. ## XIV Conclusion We have seen that one can recast obtaining a private key under LOPC in terms of distilling private states under LOCC. One finds a general class of states which are unconditionally secure. This class includes bound entangled states from which one cannot distill pure entanglement. This then enables one to use tools developed in entanglement theory to tackle privacy theory. For example, the regularized relative entropy of entanglement was found to be an upper bound on the rate of private key. Many open questions remain. The most important problem in this context is whether all entangled states have non-zero distillable key or opposite - if there are bound entangled states which cannot be distilled into private states. One can also ask about the private state cost $`K_{qc}`$ of states $`\rho _{AB}`$. I.e. what is the dimension $`d`$ of the key part of the pdit that is required to create $`\rho _{AB}`$ under LOCC? It might even be that $`K_{qc}=K_d`$, which would enable entanglement theory to have basic laws along the lines of . The question of reversibility of creating states from private states touches another ”qualitative” problem, namely how tight is the upper bound on distillable key which is the regularised relative entropy of entanglement. Exploring the wide class of private states especially in the context of the well established theory of distillation of entanglement appears to be a necessary step in order to solve the above important problems. ## XV Appendix ### XV-A Derivation of formula (101) of Sec. VI We show that a bipartite state of four subsystems $`ABA^{}B^{}`$ given as $$\rho _{ABA^{}B^{}}=\underset{iji^{}j^{}}{}\|ij_{AB}i^{}j^{}\|A_{A^{}B^{}}^{iji^{}j^{}},$$ (200) where $`A_{iji^{}j^{}}`$ are block matrices can be written as $$\rho _{ABA^{}B^{}}=\underset{iji^{}j^{}}{}\sqrt{p_{ij}p_{i^{}j^{}}}\|ij_{AB}i^{}j^{}\|[𝒰_{ij}\sqrt{\rho _E^{ij}}\sqrt{\rho _E^{i^{}j^{}}}𝒰_{i^{}j^{}}^{}]^T.$$ (201) To see this, we first write down its total purification: $$\psi _{ABA^{}B^{}E}=\underset{ij}{}\sqrt{p_{ij}}\|ij_{AB}\|\psi _{ij}_{A^{}B^{}E}.$$ (202) The states $`\psi _{A^{}B^{}E}^{ij}`$ can be written as $$\psi _{A^{}B^{}E}^{ij}=\underset{k=1}{\overset{d_{A^{}B^{}}}{}}\lambda _k^{ij}V_{ij}\|k_{A^{}B^{}}U_{ij}𝒲\|k_{A^{}B^{}}$$ (203) Here $`U_{ij}`$ is unitary transformation acting on Eve’s system, $`V_{ij}`$ is unitary transformation acting on shield $`A^{}B^{}`$ and $`𝒲`$ is some fixed embedding of $`_{A^{}B^{}}`$ into $`_E`$ (this is needed if Eve’s systems are greater than the system $`A^{}B^{}`$): $$𝒲:_{A^{}B^{}}_E,𝒲\|k_{A^{}B^{}}=\|k_E$$ (204) where $`\|k_{A^{}B^{}},k=1,\mathrm{},d_{A^{}B^{}}`$ is a fixed basis in system $`A^{}B^{}`$, while $`\|k_E,k=1,\mathrm{},d_E`$ is a fixed basis in system $`E`$. We will also need a dual operation, which is fixed projection of space $`_E`$ into $`_{A^{}B^{}}`$: $$𝒲^{}:_E_{A^{}B^{}},$$ (205) with $`𝒲^{}\|k_E=\|k_{A^{}B^{}}\text{for}k=1,\mathrm{},d_{A^{}B^{}}`$ (206) $`𝒲^{}\|k_E=0\text{for}k>d_{A^{}B^{}}`$ (207) One then finds that $$(A_{iji^{}j^{}})^T=V_{ij}^{}𝒲^{}U_{ij}^{}\sqrt{\rho _E^{ij}}\sqrt{\rho _E^{i^{}j^{}}}U_{i^{}j^{}}𝒲V_{i^{}j^{}}$$ (208) where $`T`$ is matrix transposition. One can find, the operator $`𝒰_{ij}^{}U_{ij}𝒲V_{ij}`$ maps the space $`_{A^{}B^{}}`$ exactly onto a support of $`\rho _E^{i^{}j^{}}`$ in space $`_E`$, and the dual operator $`𝒰_{ij}=V_{ij}^{}𝒲^{}U_{ij}^{}`$ maps the support of $`\rho _E^{ij}`$ back to $`_{A^{}B^{}}`$. Finally, our state is of the form $$\rho _{ABA^{}B^{}}=\underset{iji^{}j^{}}{}\sqrt{p_{ij}p_{i^{}j^{}}}\|ij_{AB}i^{}j^{}\|[𝒰_{ij}\sqrt{\rho _E^{ij}}\sqrt{\rho _E^{i^{}j^{}}}𝒰_{i^{}j^{}}^{}]^T,$$ (209) which we aimed to show. ### XV-B The proof of lemma 5, Section X-A. We prove now that the states from a family that we have introduced in eq. (146), are indeed PPT for certain range of parameters, as it is stated in lemma 5. ###### Proof: The matrix of the state (146) after partial transposition has a form $$\rho _{ABA^{}B^{}}^\mathrm{\Gamma }=\left[\begin{array}{cccc}p(\frac{\tau _1+\tau _2}{2})^\mathrm{\Gamma }& 0& 0& 0\\ 0& (\frac{1}{2}p)\tau _2^\mathrm{\Gamma }& p(\frac{\tau _1\tau _2}{2})^\mathrm{\Gamma }& 0\\ 0& p(\frac{\tau _1\tau _2}{2})^\mathrm{\Gamma }& (\frac{1}{2}p)\tau _2^\mathrm{\Gamma }& 0\\ 0& 0& 0& p(\frac{\tau _1+\tau _2}{2})^\mathrm{\Gamma }\end{array}\right].$$ (210) Since $`\tau _1`$ and $`\tau _2`$ are separable (and hence PPT), so is their mixture. Thus extreme-diagonal blocks of the above matrix are positive. It remains to check positivity of the middle block matrix. Since any block matrix of the form $$\left[\begin{array}{cc}A& B\\ B& A\end{array}\right],$$ (211) is positive if there holds $`A\|B\|`$ where $`A`$ and $`B`$ are arbitrary hermitian matrices, our question of positivity of (210) reads $$(\frac{1}{2}p)\tau _2^\mathrm{\Gamma }p\|(\frac{\tau _1\tau _2}{2})^\mathrm{\Gamma }\|$$ (212) Having $`\rho _s=\frac{1}{d^2+d}(I+V)`$ and $`\rho _a=\frac{1}{d^2d}(IV)`$ where $`V`$ swaps d-dimensional spaces and applying $`V^\mathrm{\Gamma }=dP_+`$ one easily gets that $`\tau _1^\mathrm{\Gamma }=\left({\displaystyle \frac{P_+^{}}{d^21}}\right)^k`$ (213) $`\tau _2^\mathrm{\Gamma }=\left({\displaystyle \frac{P_+^{}}{d^2+d}}+{\displaystyle \frac{(1+d)P_+}{d^2+d}}\right)^k`$ (214) where $`P_+^{}IP_+`$ is projector onto subspace orthogonal to the projector onto maximally entangled state $`P_+=\|\psi _+\psi _+\|`$. We check then the inequality $`({\displaystyle \frac{1}{2}}p)\left({\displaystyle \frac{P_+^{}}{d^2+d}}+{\displaystyle \frac{(1+d)P_+}{d^2+d}}\right)^k`$ $`{\displaystyle \frac{p}{2}}\times \left\|\left({\displaystyle \frac{P_+^{}}{d^21}}\right)^k\left({\displaystyle \frac{P_+^{}}{d^2+d}}+{\displaystyle \frac{(1+d)P_+}{d^2+d}}\right)^k\right\|`$ To solve this inequality it is useful to represent the term on LHS as a sum: $$\left(\frac{P_+^{}}{d^2+d}+\frac{(1+d)P_+}{d^2+d}\right)^k=\left(\frac{P_+^{}}{d^2+d}\right)^k+R$$ (216) where operator $`R`$ is an unnormalised state which consists of all terms coming out of $`k`$-fold tensor product of $`\left(\frac{P_+^{}}{d^2+d}+\frac{(1+d)P_+}{d^2+d}\right)`$ apart from the first term $`\left(\frac{P_+^{}}{d^2+d}\right)^k`$. It is good to note that $`R`$ has support on subspace orthogonal to $`(P_+^{})^k`$. This fact allows to omit the modulus and to get $`({\displaystyle \frac{1}{2}}p)\left[\left({\displaystyle \frac{P_+^{}}{d^2+d}}\right)^k+R\right]`$ $`{\displaystyle \frac{p}{2}}\left[(P_+^{})^k\left({\displaystyle \frac{1}{(d^21)^k}}{\displaystyle \frac{1}{(d^2+d)^k}}\right)+R\right]`$ (217) Since $`R`$ and $`(P_+^{})^k`$ are orthogonal, this inequality is equivalent to the following two inequalities $`({\displaystyle \frac{1}{2}}{\displaystyle \frac{3}{2}}p)R0`$ (218) $`({\displaystyle \frac{1}{2}}p)\left({\displaystyle \frac{P_+^{}}{d^2+d}}\right)^k{\displaystyle \frac{p}{2}}(P_+^{})^k\times `$ $`\times \left({\displaystyle \frac{1}{(d^21)^k}}{\displaystyle \frac{1}{(d^2+d)^k}}\right)`$ (219) To save first inequality one needs $`p\frac{1}{3}`$. Preserving the second one requires $$\frac{1p}{p}\left(\frac{d}{d1}\right)^k$$ (220) This however is fulfilled for any $`p(0,\frac{1}{3}]`$ if $`d`$ is taken properly large for some fixed k. Indeed, the $`k`$-th root of $`\frac{1p}{p}`$ (which converges to $`1`$ with k) can be greater than $`\frac{d}{d1}`$ (which converges to $`1`$ with d) for some large d. ### XV-C Comparison of two criteria for secure key In this section, we shall compare the joint cryptographic criterion, i.e. the requirement of (117): $$\rho _{real}^{ccq}\rho _{ideal}^{ccq}ϵ$$ (221) which includes both uniformity and security in one formula with the double condition where uniformity and security are treated separately, namely: $`\chi (\{p_i,\rho _{ij}^E\})ϵ`$ (222) $`\rho _{AB}\rho _{ideal}^{AB}ϵ`$ The connection between these two criteria for quantum cryptographical security of the state is given in the theorem below. ###### Theorem 10 For any ccq state $`\rho _{ABE}=_{ij=0}^{d1}p_{ij}\|ijij\|\rho _{ij}^E`$ and $`\rho _{ideal}=_{i=0}^{d1}\frac{1}{d}\|iiii\|\rho _E`$ where $`\rho _E=_{ij}p_{ij}\rho _{ij}^E`$, the following implications holds: $$\begin{array}{c}\rho _{AB}\rho _{ideal}^{AB}ϵ\hfill \\ \chi (\rho _{ABE})ϵ\hfill \end{array}\}\rho _{ABE}\rho _{ideal}ϵ+\sqrt{ϵ}$$ (223) $$\rho _{ABE}\rho _{ideal}ϵ\{\begin{array}{c}\chi (\rho _{ABE})4ϵ\mathrm{log}d+h(ϵ)\hfill \\ \rho _{AB}\rho _{ideal}^{AB}ϵ.\hfill \end{array}$$ (224) where $`\rho _{AB}=\mathrm{Tr}_E\rho _{ABE}`$, $`\rho _{ideal}^{AB}=\mathrm{Tr}_E\rho _{ideal}`$. ###### Remark 4 We see that the result (224) is not fully satisfactory due to the term $`\mathrm{log}d`$. However, one cannot get a better result. Indeed it is easy to construct a state, for which the Holevo function is of $`ϵ\mathrm{log}d`$ order, though the state $`\rho _{ABE}`$ is $`ϵ`$ close to some $`\rho _{ideal}`$ state. As an example may serve an appropriate extension of the isotropic state, measured in computational basis: $`\rho _{ABE}=(1ϵ)({\displaystyle \underset{i=0}{\overset{d1}{}}}{\displaystyle \frac{1}{d}}\|iiii\|_{AB})(\|0000\|)_E+`$ $`+ϵ{\displaystyle \underset{ij}{}}{\displaystyle \frac{1}{d^2d}}\|ijij\|_{AB}(\|ijij\|)_E`$ (225) where $`\sigma `$ is maximally mixed state. If we consider now $`\rho _{ideal}=(_{i=0}^{d1}\frac{1}{d}\|iiii\|_{AB})(\|0000\|)_E`$, it is easy to see that $`\rho _{ABE}\rho _{ideal}=2ϵ`$. However the value of the Holevo function equals $`h(ϵ)+ϵ\mathrm{log}(d^2d)`$. ###### Remark 5 The main difficulty in the proof of the above theorem is to get the term $`\mathrm{log}d`$ ($`d\times d`$ is size of $`AB`$ system) rather than $`\mathrm{log}d_{ABA^{}B^{}}`$. The latter one would be obtained directly from Fannes type continuities. However to get $`\mathrm{log}d`$ we have to apply tricks based on twisting. It is quite convenient not to have Eve’s dimension in equivalence formula. This is because Eve’s dimension depends on the protocol that lead to the key (more specifically, it depends on the amount of communication). In contrast, dimension of Alice and Bob system is only the number of bits of obtained key. Thus, our equivalence is independent of the protocol. ###### Proof: For the first part of the theorem 10 we assume that $`\chi (\{p_{ij},\rho _{ij}^E\})ϵ`$ which by proposition 7 in Sec. XV-D of Appendix means that we have: $$\underset{i,j=0}{\overset{d1}{}}p_{ij}\|ijij\|_{AB}\rho _{ij}^E\underset{i,j=0}{\overset{d1}{}}p_{ij}\|ijij\|_{AB}\rho _E\sqrt{ϵ},$$ (226) with $`\rho _E=_{i,j=0}^{d1}p_{ij}\rho _{ij}^E`$. Moreover by second assumption that $$\|\|\underset{i,j=0}{\overset{d1}{}}p_{ij}\|ijij\|\underset{i}{\overset{d1}{}}\frac{1}{d}\|iiii\|)\|\|ϵ$$ (227) one gets $`{\displaystyle \underset{i,j=0}{\overset{d1}{}}}p_{ij}\|ijij\|_{AB}\rho _E{\displaystyle \underset{i=0}{\overset{d1}{}}}{\displaystyle \frac{1}{d}}\|iiii\|_{AB}\rho _Eϵ.`$ (228) Using triangle inequality, and Eqs. (226) and (228) one obtains the $$\rho _{ABE}\rho _{ideal}ϵ+\sqrt{ϵ}.$$ (229) The proof of the second part of the theorem 10 is a bit more involved. Of course, it is immediate that due to monotonicity of trace norm under partial trace, from $`\rho _{ABE}\rho _{ideal}ϵ`$ it follows $`\rho _{AB}\rho _{ideal}^{AB}ϵ`$. The non-obvious task is to bound also $`\chi `$. So, we assume that $$\rho _{ABE}\rho _{ideal}ϵ,$$ (230) By equality of norm and fidelity condition (11), there holds $$F(\rho _{ABE},\rho _{ideal})1\frac{1}{2}ϵ$$ (231) By definition of fidelity, there are pure states $`\psi `$ and $`\varphi `$ (purifications of $`\rho _{ABE}`$ and $`\rho _{ideal}`$ respectively), such that $`F(\psi ,\varphi )=F(\rho _{ABE},\rho _{ideal})`$. Without loss of generality we can consider the system which purifies both states to be bipartite. We will call it $`A^{}B^{}`$. Now let us perform twisting operation on the $`ABA^{}B^{}`$ parts of the pure states $`\psi `$ and $`\varphi `$, which in the case of state $`\rho _{ideal}`$ transforms $`AB`$ subsystem of into maximally entangled state - $`P_d^+`$, (we can choose such twisting because by the theorem 2 purification of an ideal state is some pdit state). I.e. after such twisting, pdit will become a basic pdit (7) which is product with $`A^{}B^{}`$ subsystem. Since unitary transformation and tracing out can only increase fidelity, then applying again (11) we have that subsystem $`AB`$ of $`\rho _{ABE}`$ is close to a singlet state in norm: $$\rho _{AB}P_d^+2\sqrt{ϵ}.$$ (232) Using Fannes inequality (see eq. 12 in Sec. II-A) we get $$S(\rho _{AB})4\sqrt{ϵ}\mathrm{log}d+h(2\sqrt{ϵ}).$$ (233) Since the total state of systems $`ABA^{}B^{}E`$ is pure, we get that $`S(\rho _{A^{}B^{}E})=S(\rho _{AB})`$ hence $$S(\rho _{A^{}B^{}E})4\sqrt{ϵ}\mathrm{log}d+h(2\sqrt{ϵ}).$$ (234) Now, note that the state of the system $`A^{}B^{}E`$ has the form $$\rho _{A^{}B^{}E}=\underset{i,j=0}{\overset{d1}{}}p_{ij}\rho _{ij}^{A^{}B^{}E},$$ (235) where the state $`\rho _{ij}^{A^{}B^{}E}`$ denotes state of $`A^{}B^{}E`$ system after twisting and given that $`AB`$ subsystem is in state $`\|ijij\|`$ (i.e. if after twisting one measure the system $`AB`$ in basis $`\|ij`$ the system $`A^{}B^{}E`$ would collapse to $`\rho _{ij}^{A^{}B^{}E}`$). By definition of Holevo function there holds: $$\chi (\{p_{ij},\rho _{ij}^{A^{}B^{}E}\})S(\rho _{A^{}B^{}E}).$$ (236) The question is how the Holevo function of the $`\{p_{ij},\rho _{ij}^{A^{}B^{}E}\}`$ ensemble is related to Holevo function of $`\{p_{ij},\rho _{ij}^E\}`$ which we would like to bound from above. It is crucial, that by theorem 1 twisting operation does not affect the ccq state which comes out of the measurement of $`AB`$ in control basis of twisting. In other words, the ensemble $`\{p_i,\rho _{ij}^E\}`$ does not change under twisting, so that $`\rho _{ij}^E=\mathrm{Tr}_{A^{}B^{}}\rho _{ij}^{A^{}B^{}E}`$. It is easy now to compare the functions $`\chi (\{p_{ij},\rho _{ij}^E\})`$ and $`\chi (\{p_{ij},\rho _{ij}^{A^{}B^{}E}\})`$: $$\chi (\{p_{ij},\rho _{ij}^E\})\chi (\{p_{ij},\rho _{ij}^{A^{}B^{}E}\})$$ (237) This is due to the fact, that each state $`\rho _{ij}^E`$ can be obtained from $`\rho _{ij}^{A^{}B^{}E}`$ by tracing out $`A^{}B^{}`$ subsystem. However tracing out can only decrease Holevo function, because this function is equal to the average relative entropy: $$\chi (\{p_k,\rho _k\})=\underset{k}{}p_kS(\rho _k\|\underset{k}{}p_k\rho _k).$$ (238) Summing up the chain of inequalities (234), (236) and (237) one gets $$\chi (\{p_{ij},\rho _{ij}^E\})4\sqrt{ϵ}\mathrm{log}d+h(2\sqrt{ϵ}).$$ (239) which is a desired security condition - bound on the Holevo function of the ansamble $`\{p_{ij},\rho _{ij}^E\}`$. ### XV-D Useful inequalities relating security conditions In this section we collect relations between different security conditions for ccq states. Some of these relations have been studied in . Since we will not deal with uniformity, but solely with security, it is convenient to use single index $`k`$ in place of $`ij`$. We thus consider ccq state (which could be actually called cq state) $$\rho =\underset{k=0}{\overset{d^21}{}}p_k\|kk\|\rho _k$$ (240) Basing on this fact, we can state another lemma establishing some equivalences: ###### Lemma 8 For any state (240) and any positive real $`ϵ\frac{1}{2}`$, the following implications hold 1. $`{\displaystyle \underset{k}{}}p_k\|kk\|\rho _k({\displaystyle \underset{j}{}}p_j\|jj\|)\rho ϵ`$ (241) $`{\displaystyle \underset{k}{}}p_kF(\rho _k,\rho )1{\displaystyle \frac{1}{2}}ϵ`$ 2. $$\underset{k}{}p_kF(\rho _k,\rho )1ϵ\underset{k}{}p_k\rho _k\rho 8ϵ$$ 3. $`{\displaystyle \underset{k}{}}p_k\rho _k\rho ϵ`$ (242) $`{\displaystyle \underset{k}{}}p_k\|kk\|\rho _k({\displaystyle \underset{j}{}}p_j\|jj\|)\rho ϵ.`$ Here $`\rho =_kp_k\rho _k`$. ###### Proof: The first thesis follows from the mentioned equivalence of norm and fidelity and definition of fidelity. Namely one can make use of lemma 1, so that if (241) holds, the fidelity $`F(_kp_k\|kk\|\rho _k,(_jp_j\|jj\|)\rho )`$ is no less than $`1\frac{1}{2}ϵ`$. However it is equal to average fidelity $`_kp_kF(\rho _k,\rho )`$. Indeed, $$F(\underset{k}{}p_k\|kk\|\rho _k,(\underset{j}{}p_j\|jj\|)\rho )=\mathrm{Tr}((\underset{k}{}p_k\|kk\|\rho _k)^{\frac{1}{2}}\underset{j}{}p_j\|jj\|\rho ..(\underset{l}{}p_l\|ll\|\rho _l)^{\frac{1}{2}})^{\frac{1}{2}}$$ (243) Now by orthogonality of vectors $`\|k`$ one has $$\sqrt{\underset{k}{}p_k\|kk\|\rho _k}=\underset{k}{}\sqrt{p_k}\|kk\|\sqrt{\rho _k}.$$ (244) Multiplying now the $`(_jp_j\|jj\|)\rho `$ matrix by the above from left-hand-side and right-hand-side one gets $$\underset{k}{}p_k^2\|kk\|\sqrt{\rho _k}\rho \sqrt{\rho _k}.$$ (245) This immediately gives the above formula equal to $$\underset{k}{}p_k\mathrm{Tr}\sqrt{\sqrt{\rho _k}\rho \sqrt{\rho _k}}$$ (246) which is just average fidelity from (2). The second thesis of this lemma ( Eq. (2)) is again a consequence of (11). If applied to each pair $`\rho _k`$, $`\rho `$, and averaged over probabilities of $`p_k`$ gives that $$\underset{k}{}p_k\sqrt{1\frac{1}{4}\rho \rho _k^2}1ϵ$$ (247) which is equivalent to $$\underset{k}{}p_k\left(1\sqrt{1\rho \rho _k^2/4}\right)ϵ.$$ (248) Now by the fact that $$1\sqrt{1\frac{1}{4}\rho \rho _k^2}$$ (249) is a convex function of $`\rho \rho _k`$ on interval $`(0,2)`$ we get $$1\sqrt{1\frac{1}{4}[\underset{k}{}p_k\rho \rho _k]^2}ϵ$$ (250) This however reads for $`0<ϵ<1`$ $$\underset{k}{}p_k\rho \rho _k8ϵ.$$ (251) Since $`\rho \rho _k2`$ one has, that for $`ϵ1`$ the above inequality is also valid, which completes the proof of the second thesis of lemma 8. The last implication (Eq. (242)) is a consequence of triangle inequality, which completes the lemma. Let us notice, that this lemma establishes a kind of equivalance of security conditions, namely: $`{\displaystyle \underset{k}{}}p_k\|kk\|\rho _k({\displaystyle \underset{j}{}}p_j\|jj\|)\rho ϵ`$ (252) $`{\displaystyle \underset{k}{}}p_k\rho \rho _k4ϵ`$ $`{\displaystyle \underset{k}{}}p_k\|kk\|\rho _k({\displaystyle \underset{j}{}}p_j\|jj\|)\rho 4ϵ`$ We can show now links between the above conditions on ccq state and Holevo function $`\chi `$ of this state, i.e. of an ansamble $`\{p_k,\rho _k\}`$ which we shall write $`\chi (\rho _{ccq})`$. ###### Lemma 9 For any ccq state $`\rho _{ccq}`$ (240) there holds: $`\chi (\rho _{ccq})ϵ{\displaystyle \underset{k}{}}p_k\rho _k\rho \sqrt{2ϵ}`$ (253) $`{\displaystyle \underset{k}{}}p_k\rho _k\rho ϵ\chi (\rho _{ccq})3ϵ\mathrm{log}d+\mathrm{max}(h(ϵ),2ϵ)`$ where $`_kp_k\rho _k=\rho `$, which acts on Hilbert space $`=𝒞^d`$, and $`h(ϵ)=ϵ\mathrm{log}ϵ(1ϵ)\mathrm{log}(1ϵ)`$ is binary entropy. ###### Proof: For the first statement of this lemma, let us notice that $`\chi (\rho _{ccq})=S(\rho )_kp_kS(\rho _k)`$ is just equal to average relative entropy distance $`_kp_kS(\rho _k\|\rho )`$. Thus, by assumption we have $$\chi (\rho _{ccq})=\underset{k}{}p_kS(\rho _k\|\rho )ϵ.$$ (254) Now we can make use of the inequality : $$\frac{1}{2}\rho \rho _k^2S(\rho _k\|\rho )$$ (255) which after averaging over probabilities and by concavity of square root gives $$\underset{k}{}p_k\rho \rho _k\sqrt{2\underset{k}{}p_kS(\rho _k\|\rho )}.$$ (256) Applying now bound (254) we obtain $$\underset{k}{}p_k\rho \rho _k\sqrt{2ϵ}$$ (257) which completes first thesis of this lemma. To prove the second statement of the lemma we use the Fannes inequality (see eq. 12 in Sec. II-A). Namely, for $`\rho \rho _k1`$ there holds: $$\|S(\rho )S(\rho _k)\|2\rho \rho _k\mathrm{log}d+h(\rho \rho _k).$$ (258) Let $`G=\{k:\rho \rho _k1\}`$ and $`B=\{k:\rho \rho _k>1\}`$, and denote $`_{kG}p_k[S(\rho )S(\rho _k)]\chi _G(\rho _{ccq})`$, and $`\chi _B(\rho _{ccq})`$ analogously. We then have, that $$\chi (\rho _{ccq})=\chi _G(\rho _{ccq})+\chi _B(\rho _{ccq}).$$ (259) We will give now the bounds for $`\chi _G(\rho _{ccq})`$ and $`\chi _B(\rho _{ccq})`$ respectively. The first quantity is directly bounded by the analogous sum over LHS of the Fannes inequality: $$\chi _G(\rho _{ccq})\underset{kG}{}p_k\left(2\rho \rho _k\mathrm{log}d+h(\rho \rho _k)\right).$$ (260) ¿From assumption $$\underset{k}{}p_k\rho \rho _kϵ$$ (261) it follows that $`_{kG}p_k\rho \rho _kϵ`$. Using this, and adding non-negative terms $`_{kB}p_kh(\rho \rho _k)`$, we get: $$\chi _G(\rho _{ccq})2ϵ\mathrm{log}d+\underset{k}{}p_kh(\rho \rho _k).$$ (262) Now by concavity of binary entropy one gets $$\chi _G(\rho _{ccq})2ϵ\mathrm{log}d+h(\underset{k}{}p_k\rho \rho _k),$$ (263) Were the entropy increasing on $`[0,\mathrm{}]`$ interval, one could use directly the assumption that $`_kp_k\rho \rho _kϵ`$, and bound $`h(_kp_k\rho \rho _k)`$ by $`h(ϵ)`$. Since it is the case only for $`ϵ[0,\frac{1}{2}]`$, we have to end up with more ugly, but nonetheless useful expression. Namely on the interval $`(\frac{1}{2},\mathrm{}]`$ where the entropy becomes decreasing, it is bounded by 1, and hence not greater than $`2ϵ`$ for $`ϵ(\frac{1}{2},\mathrm{}]`$. Thus finally one gets $$\chi _G(\rho _{ccq})2ϵ\mathrm{log}d+\mathrm{max}(h(ϵ),2ϵ),$$ (264) We turn now to give the bound for $`\chi _B(\rho _{ccq})`$. For the latter we use the fact, that the Holevo quantity is bounded from above by $`\mathrm{log}d`$, which gives: $$\chi _B(\rho _{ccq})\underset{kB}{}p_k\mathrm{log}d.$$ (265) To bound the last inequality we observe, that by definition of the set $`B`$ we have $`_{kB}p_k\rho \rho _k_{kB}p_k`$. Then, again by assumption (261), we have $$\chi _B(\rho _{ccq})ϵ\mathrm{log}d.$$ (266) Collecting inequality (264) and the above one, we arrive at the formula $$\chi (\rho _{ccq})3ϵ\mathrm{log}d+\mathrm{max}(h(ϵ),2ϵ),$$ (267) which ends the proof of the lemma. The lemmas above allow to prove the following proposition ###### Proposition 7 For any ccq state (240) the following holds: $`\chi (\{p_k,\rho _k\})ϵ`$ $`{\displaystyle \underset{k}{}}p_k\|kk\|\rho _k({\displaystyle \underset{j}{}}p_j\|jj\|)\rho \sqrt{2ϵ}`$ (268) where $`\rho =_kp_k\rho _k`$. ###### Proof: Assuming that Holevo function is smaller than $`ϵ`$, we get by lemma 9 that $`_kp_k\rho _k\rho \sqrt{2ϵ}`$. This however implies by lemma (8) that $`_kp_k\|kk\|\rho _k(_jp_j\|jj\|)\rho `$ is also not greater than $`\sqrt{2ϵ}`$, which completes proof of the proposition. ### XV-E Properties of pbits We shall give here detailed proof of the lemma 3, Section IV-A. ###### Proof: Log-negativity (cf. ) is defined as $`E_N(\rho )=\mathrm{log}(\rho ^\mathrm{\Gamma })`$. It is easy to see, that after partial transposition on $`BB^{}`$ subsystem, the pbit $`\gamma `$ in $`X`$-form changes into $$\gamma _{ABA^{}B^{}}^\mathrm{\Gamma }=\frac{1}{2}\left[\begin{array}{cccc}\sqrt{XX^{}}^\mathrm{\Gamma }& 0& 0& 0\\ 0& 0& X^\mathrm{\Gamma }& 0\\ 0& (X^{})^\mathrm{\Gamma }& 0& 0\\ 0& 0& 0& \sqrt{X^{}X}^\mathrm{\Gamma }\end{array}\right].$$ (269) We have $$\gamma ^\mathrm{\Gamma }=\frac{1}{2}([\sqrt{XX^{}}]^\mathrm{\Gamma }+[\sqrt{X^{}X}]^\mathrm{\Gamma }+A),$$ (270) where $$A=\left[\begin{array}{cc}0& (X)^\mathrm{\Gamma }\\ (X^{})^\mathrm{\Gamma }& 0\end{array}\right].$$ (271) By assumption, the operators $`[XX^{}]^\mathrm{\Gamma }`$ and $`[X^{}X]^\mathrm{\Gamma }`$ are positive, so that $$[\sqrt{XX^{}}]^\mathrm{\Gamma }+[\sqrt{X^{}X}]^\mathrm{\Gamma }=\mathrm{Tr}(\sqrt{XX^{}}+\sqrt{X^{}X})^\mathrm{\Gamma }=2\mathrm{T}\mathrm{r}\gamma ^\mathrm{\Gamma }=2.$$ (272) The last equality comes from the fact that $`\mathrm{\Gamma }`$ preserves trace. To evaluate norm of $`A`$, we note that due to unitary invariance of trace norm we have $`A=\sigma _xI_{A^{}B^{}}A`$. Consequently $$A=X^\mathrm{\Gamma }+(X^{})^\mathrm{\Gamma }=2X^\mathrm{\Gamma }.$$ (273) The last equality follows form the fact that $`\mathrm{\Gamma }`$ commutes with Hermitian conjugation, and trace norm is invariant under Hermitian conjugation $`X=X^{}`$. Thus we get $$E_N(\gamma )=\mathrm{log}(1+X^\mathrm{\Gamma }),$$ (274) which proves the lemma. ### XV-F Relative entropy of entanglement and pdits We give now the proof of the theorem 66, Section XV-F. ###### Proof: If one consider the thesis of theorem 3 for the state $`\rho =\gamma _{ABA^{}B^{}}^n`$, it follows that $$E_r(\rho )\mathrm{log}d^n+\frac{1}{d^n}\underset{k=0}{\overset{d^n1}{}}E_r(\sigma _k),$$ (275) with $`k`$ being the multiindex $`k=(i_1,\mathrm{},i_n)`$ with $`i_l\{0,\mathrm{},d1\}`$ for $`l\{1,\mathrm{},n\}`$ and $`\sigma _k=\rho _{i_1}\mathrm{}\rho _{i_n}`$. Dividing both sides by $`n`$ we obtain $$\frac{1}{n}E_r(\gamma _{ABA^{}B^{}}^n)\mathrm{log}d+\frac{1}{nd^n}\underset{k=0}{\overset{d^n1}{}}E_r(\sigma _k),$$ (276) The left-hand-side of this inequality approaches $`E_r^{\mathrm{}}(\gamma )`$ with $`n`$. What has to be shown is that $$\underset{n\mathrm{}}{lim}\frac{1}{nd^n}\underset{k=0}{\overset{d^n1}{}}E_r(\sigma _k)\underset{i=0}{\overset{d1}{}}\frac{1}{d}E_r^{\mathrm{}}(\rho _i),$$ (277) with $`\rho _i`$ denoting the conditional states on $`A^{}B^{}`$ subsystem. Let us first observe that $`E_r(\sigma _k)=E_r(\sigma _k^{})`$ for any $`k`$ and $`k^{}`$ which are of the same type, i.e. which has the same numbers of occurrence of symbols from set $`\{0,\mathrm{},d1\}`$. This is because $`\sigma _k`$ and $`\sigma _k^{}`$ differ by local reversible transformation which does not change the entanglement. Moreover, as we will see, one can consider only those $`\sigma _k`$ for which $`k`$ is $`\delta `$-strongly typical i.e. such, that for some fixed $`\delta >0`$ there holds : $$_{a\{0,\mathrm{},d1\}}\|\frac{a(k)}{n}\frac{1}{d}\|<\delta ,$$ (278) where $`a(k)`$ denotes frequency of symbol $`a`$ in sequence $`k`$. The set of such $`k`$ of length $`n`$ we will denote as $`ST_\delta ^n`$. It is known, that the strongly typical set carries almost whole probability mass for large $`n`$, that is for any $`\delta >0`$ and any $`ϵ>0`$ there exists $`n_0`$ such that for all $`nn_0`$: $$P^n(ST_\delta ^n)>1ϵ.$$ (279) Now, since we deal here with homogeneous distribution, we can say, that the probability of the set of events is directly related to the power of this set. Namely we have: $$\|ST_\delta ^n\|/d^n>1ϵ,$$ (280) which gives $`\|ST_\delta ^n\|>(1ϵ)d^n`$. We can rewrite now the term of the LHS of (277) as follows: $$\frac{1}{nd^n}\underset{k=0}{\overset{d^n1}{}}E_r(\sigma _k)=\frac{1}{nd^n}\left(\underset{kST_\delta ^n}{}E_r(\sigma _k)+\underset{kST_\delta ^n}{}E_r(\sigma _k)\right).$$ (281) We can get rid of the second term of the RHS of the above equality, because we can bound from above for each $`k`$ the term $`E_r(\sigma _k)`$ by $`n\mathrm{log}d`$, which gives: $$\frac{1}{nd^n}\underset{k=0}{\overset{d^n1}{}}E_r(\sigma _k)\frac{1}{nd^n}\left(\underset{kST_{\delta ^n}}{}E_r(\sigma _k)+n\mathrm{log}d\underset{kST_\delta ^n}{}\right)\frac{1}{nd^n}\underset{kST_\delta ^n}{}E_r(\sigma _k)+ϵ\mathrm{log}d$$ (282) Now for each $`kST_\delta ^n`$ we have $$E_r(\sigma _k)=E_r(\rho _0^{\stackrel{~}{m}_0}\rho _1^{\stackrel{~}{m}_1}\mathrm{}\rho _{d1}^{\stackrel{~}{m}_{d1}})$$ (283) with $`\stackrel{~}{m}_i=i(k)`$ with $`i`$ in place of $`a`$ in (278). By subadditivity of $`E_r`$ one has $$E_r(\sigma _k)\underset{l=0}{\overset{d1}{}}E_r(\rho _l^{\stackrel{~}{m}_l}).$$ (284) Note, that $`\rho _l`$ stands here for the state on shield part of one copy of $`\gamma _{ABA^{}B^{}}`$. Applying this inequality for each $`k`$ in $`ST_\delta ^n`$ and taking maximum of LHS of the above inequality over $`k`$, we have a bound: $$\frac{1}{nd^n}\underset{k=0}{\overset{d^n1}{}}E_r(\sigma _k)\underset{l=0}{\overset{d1}{}}\frac{1}{n}E_r(\rho _l^{m_l})+ϵ\mathrm{log}d,$$ (285) where $`m_l`$ are the coefficients of the decomposition of some $`\sigma _k`$ into $`\rho ^{m_l}`$, which yields the maximal value of $`E_r(\sigma _k)`$ over all strongly typical $`k`$. Now, we can rewrite the RHS of the above inequality as: $$\underset{l=0}{\overset{d1}{}}\frac{m_l}{n}\frac{1}{m_l}E_r(\rho _l^{m_l})+ϵ\mathrm{log}d(\frac{1}{d}+\delta )\underset{l=0}{\overset{d1}{}}\frac{1}{m_l}E_r(\rho _l^{m_l})+ϵ\mathrm{log}d$$ (286) where the last inequality holds for sufficiently high $`n`$ by assumption of strong typicality. Thus, for every $`\delta `$ and $`ϵ`$ and sufficiently large $`n`$ there holds: $$\frac{1}{nd^n}\underset{k=0}{\overset{d^n1}{}}E_r(\sigma _k)(\frac{1}{d}+\delta )\underset{l=0}{\overset{d1}{}}\frac{1}{m_l}E_r(\rho _l^{m_l})+ϵ\mathrm{log}d,$$ (287) One then sees, that the RHS approaches $$(\frac{1}{d}+\delta )\underset{l=0}{\overset{d1}{}}E_r^{\mathrm{}}(\rho _l)+ϵ\mathrm{log}d$$ (288) in limit of large $`n`$. Indeed, since for every $`l`$ the limit $`E_r^{\mathrm{}}(\rho _l)`$ exists, the subsequence $`\frac{1}{m_l}E_r(\rho _l^{m_l})`$ approaches this limit. Taking now infimum over $`\delta `$ and $`ϵ`$ we prove the inequality (277). This proves the theorem 66. ### XV-G Approximate pbits We give here the proof of lemma 68 ###### Proof: Assume first, that $`\mathrm{Tr}\rho _{AB}P_+>1ϵ`$. Since the elements $`a_{ijkl}`$ are real, by hermicity of the state we have $$\mathrm{Tr}\rho _{AB}P_+=\frac{1}{2}(a_{0000}+a_{1111}+2Re(a_{0011}))$$ (289) This is however less than or equal to $`\frac{1}{2}(1+2a_{0011})`$, which is in turn greater than $`1ϵ`$, and the assertion follows. For the second part of the lemma, assume that $`a_{0011}>\frac{1}{2}ϵ`$. We then have $$\mathrm{Tr}\rho _{AB}P_+>\frac{1}{2}(a_{0000}+a_{1111}+12ϵ).$$ We now bound the sum of $`a_{0000}`$ and $`a_{1111}`$. By positivity of the state, we have that $`\sqrt{a_{0000}a_{1111}}>\|a_{0011}\|`$. Now, by arithmetic-geometric mean inequality, we have that $`a_{0000}+a_{1111}2\sqrt{a_{0000}a_{1111}}`$ which gives the proof. ### XV-H Relative entropy bound ###### Proof: (of Lemma 7, Section XI) Let us first show, that $$\mathrm{Tr}P_d^+\sigma _{AB}\frac{1}{d}$$ (290) for any $`\sigma _{AB}T`$. We first show this for $`\sigma _{AB}`$ ”derived” from some pure product states $`\|\psi \psi \|`$: $$\sigma _{AB}=\mathrm{Tr}_{A^{}B^{}}U^{}\|\psi \psi \|U.$$ (291) Because $`\psi `$ is product, it can be written as $$\psi =(a_i\|i_A\|\psi _i)(b_i\|i_B\|\varphi _i)$$ (292) with $`a_i,b_i`$ normalized and $`\|i_A,\|i_B,\|\psi _i,\|\varphi _i`$ on subsystem $`A,B,A^{},B^{}`$ respectively. Now the condition that the reduced $`AB`$ state has overlap with $`P_d^+`$ no greater than $`1/d`$ is $$\underset{ij}{}a_ib_ia_j^{}b_j^{}x_i\|x_j1$$ (293) where $`x_k`$ are arbitrary vectors of norm one arising from the action of $`U`$ on $`\psi _i`$ and $`\varphi _i`$. Since the $`x_k`$ are arbitrary they can incorporate the phases of $`a_i,b_i`$ so that we require now $`_{ij}\sqrt{p_iq_ip_jq_j}x_i\|x_j1`$. where $`p_i`$ and $`q_i`$ are probabilities. Now, the right hand side will not decrease if we assume $`x_i\|x_j=1`$ so we require $`[_i\sqrt{p_iq_i}]^21`$ which is satisfied by any probability distribution, which gives the proof of (290) for special $`\sigma _{AB}`$. To show the inequality is true in general we find that $`\mathrm{Tr}P_d^+\mathrm{Tr}_{A^{}B^{}}U^{}{\displaystyle \underset{k}{}}p_k\|\psi _k\psi _k\|U=`$ $`{\displaystyle \underset{k}{}}p_k\mathrm{Tr}P_d^+\mathrm{Tr}_{A^{}B^{}}U^{}\|\psi _k\psi _k\|U.`$ (294) Thus if (290) holds for $`\sigma _{AB}`$ derived from pure (product) state, by averaging over probabilities, we will have (290) for an arbitrary $`\sigma _{AB}`$ from the set $`T`$. Now by concavity of logarithm, we have for any states $`\rho `$ and $`\sigma `$: $`S(\rho \|\|\sigma )=S(\rho )\mathrm{Tr}(\rho \mathrm{log}\sigma )`$ $`S(\rho )\mathrm{log}(\mathrm{Tr}\rho \sigma )`$ (295) Applying inequality (290) we have that $$\mathrm{log}(\mathrm{Tr}\rho \sigma )\mathrm{log}d.$$ (296) Now by (295) we have that $$S(P_d^+\|\|\sigma _{AB})\mathrm{log}d,$$ (297) which is a desired bound. ### Acknowledgements We would like to thank Matthias Christandl, Ryszard Horodecki, Andrzej Szepietowski and Andreas Winter for helpful discussion. We also thank Geir Ove Myhr for fixing numerous bugs which appeared in the first version of the manuscript. This work is supported by the Polish Ministry of Scientific Research and Information Technology under the (solicited) grant No. PBZ-MIN-008/P03/2003, EU grants RESQ (IST-2001-37559), QUPRODIS (IST-2001-38877) and PROSECCO (IST-2001-39227). JO acknowledges the Cambridge-MIT Institute. We acknowledge hospitality of the Isaac Newton Institute for Mathematical Sciences during the QIS programme where the part of this work was done.
warning/0506/hep-ex0506053.html	ar5iv	text	# Absolute Branching Fraction Measurements of Exclusive 𝐷⁺ Semileptonic Decays CLEO Collaboration ## Abstract Using data collected at the $`\psi (3770)`$ resonance with the CLEO-c detector at the Cornell $`e^+e^{}`$ storage ring, we present improved measurements of the absolute branching fractions of $`D^+`$ decays to $`\overline{K}^0e^+\nu _e`$, $`\pi ^0e^+\nu _e`$, $`\overline{K}^0e^+\nu _e`$, and $`\rho ^0e^+\nu _e`$, and the first observation and absolute branching fraction measurement of $`D^+\omega e^+\nu _e`$. We also report the most precise tests to date of isospin invariance in semileptonic $`D^0`$ and $`D^+`$ decays. preprint: CLNS 05-1915preprint: CLEO 05-07 The quark mixing parameters are fundamental constants of the Standard Model (SM) of particle physics. They determine the nine weak-current quark coupling elements of the Cabibbo-Kobayashi-Maskawa (CKM) matrix ckm . The extraction of the quark couplings is difficult because quarks are bound inside hadrons by the strong interaction. Semileptonic decays are the preferred way to determine the CKM matrix elements as the strong interaction binding effects are confined to the hadronic current. They are parameterized by form factors that are calculable, for example, by lattice quantum chromodynamics (LQCD) and QCD sum rules. Nevertheless, form factor uncertainties dominate the precision with which the CKM matrix elements can be determined VubVcb . In charm quark decays, however, couplings $`V_{cs}`$ and $`V_{cd}`$ are tightly constrained by the unitarity of the CKM matrix. Therefore, measurements of charm semileptonic decay rates and form factors rigorously test theoretical predictions. We report herein measurements with the first CLEO-c data cleoc of the absolute branching fractions of $`D^+`$ decays to $`\overline{K}^0e^+\nu _e`$, $`\pi ^0e^+\nu _e`$, $`\overline{K}^0e^+\nu _e`$, and $`\rho ^0e^+\nu _e`$, and the first observation and absolute branching fraction measurement of $`D^+\omega e^+\nu _e`$. (Throughout this Letter charge-conjugate modes are implied.) We combine these results with the measurements of $`D^0`$ semileptonic branching fractions reported in Ref. cleoc-neutral\_semilep , which use the same data and analysis technique, and test isospin invariance of the hadronic current in semileptonic decays. The data were collected by the CLEO-c detector at the $`\psi (3770)`$ resonance, about 40 MeV above the $`D\overline{D}`$ pair production threshold. A description of the CLEO-c detector is provided in Ref. cleoc-neutral\_semilep and references therein. The data sample consists of an integrated luminosity of 55.8 pb<sup>-1</sup> and includes about 0.16 million $`D^+D^{}`$ events. The technique for these measurements was first applied by the Mark III collaboration MkIII at SPEAR. Candidate events are selected by reconstructing a $`D^{}`$, called a tag, in the following six hadronic final states: $`K_S^0\pi ^{}`$, $`K^+\pi ^{}\pi ^{}`$, $`K_S^0\pi ^{}\pi ^0`$, $`K^+\pi ^{}\pi ^{}\pi ^0`$, $`K_S^0\pi ^{}\pi ^{}\pi ^+`$, and $`K^{}K^+\pi ^{}`$. The absolute branching fractions of $`D^+`$ semileptonic decays are then measured by their reconstruction in the system recoiling from the tag. Tagging a $`D^{}`$ meson in a $`\psi (3770)`$ decay provides a $`D^+`$ with known four-momentum, allowing a semileptonic decay to be reconstructed with no kinematic ambiguity, even though the neutrino is undetected. Tagged events are selected based on two variables: $`\mathrm{\Delta }EE_DE_{\mathrm{beam}}`$, the difference between the energy of the $`D^{}`$ tag candidate ($`E_D`$) and the beam energy ($`E_{\mathrm{beam}}`$), and the beam-constrained mass $`M_{\mathrm{bc}}\sqrt{E_{\mathrm{beam}}^2/c^4\|\stackrel{}{p}_D\|^2/c^2}`$, where $`\stackrel{}{p}_D`$ is the measured momentum of the $`D^{}`$ candidate. Selection criteria for tracks, $`\pi ^0`$ and $`K_S^0`$ candidates for tags are described in Ref. cleoc-Dtagging . If multiple candidates are present in the same tag mode, one candidate per tag charge is chosen using $`\mathrm{\Delta }E`$. The yields of the six tag modes are obtained from fits to the $`M_{\mathrm{bc}}`$ distributions. The data sample comprises approximately 32,000 charged tags (Table 1). After a tag is identified, we search for a positron and a set of hadrons recoiling against the tag. (Muons are not used as $`D`$ semileptonic decays at the $`\psi (3770)`$ produce low momentum leptons for which the CLEO-c muon identification is not efficient.) Positron candidates, selected with criteria described in Ref. cleoc-neutral\_semilep , are required to have momentum of at least 200 MeV/$`c`$ and to satisfy $`\|\mathrm{cos}\theta \|`$ $`<`$ 0.90, where $`\theta `$ is the angle between the positron direction and the beam axis. The efficiency for positron identification rises from about $`50\%`$ at 200 MeV/$`c`$ to 95% just above 300 MeV/$`c`$ and is roughly constant thereafter. The rates for misidentifying charged pions and kaons as positrons averaged over the momentum range are approximately 0.1%. Bremsstrahlung photons are recovered by the procedure described in Ref. cleoc-neutral\_semilep . Hadronic tracks must have momenta above 50 MeV/$`c`$ and $`\|\mathrm{cos}\theta \|<0.93`$. Identification of hadrons is based on measurements of specific ionization ($`dE/dx`$) in the main drift chamber and information from the Ring Imaging Cherenkov Detector (RICH). Pion and kaon candidates are required to have $`dE/dx`$ measurements within three standard deviations (3.0$`\sigma `$) of the expected value. For tracks with momenta greater than 700 MeV/$`c`$, RICH information, if available, is combined with $`dE/dx`$. The efficiencies ($`95`$% or higher) and misidentification rates (a few per cent) are determined with charged pion and kaon samples from hadronic $`D`$ decays. We form $`\pi ^0`$ candidates from pairs of photons, each having an energy of at least 30 MeV, and require that the invariant mass of the pair be within 3.0$`\sigma `$ ($`\sigma 6`$ $`\mathrm{MeV}/c^2`$) of the known $`\pi ^0`$ mass. A mass constraint is imposed when $`\pi ^0`$ candidates are used in further reconstruction. The $`K_S^0`$ candidates are formed from pairs of oppositely-charged and vertex-constrained tracks having an invariant mass within 12 MeV/$`c^2`$ $`(4.5\sigma )`$ of the known $`K_S^0`$ mass. We form a $`\overline{K}^0`$ ($`\rho ^0`$) candidate from $`K^{}`$ and $`\pi ^+`$ ($`\pi ^{}`$ and $`\pi ^+`$) candidates and require an invariant mass within 100 MeV/$`c^2`$ (150 MeV/$`c^2`$) of its mean value. The reconstruction of $`\omega \pi ^+\pi ^{}\pi ^0`$ candidates is achieved by combining three pions, requiring an invariant mass within 20 MeV/$`c^2`$ of the known mass, and demanding that the charged pions do not satisfy interpretation as a $`K_S^0`$. The tag and the semileptonic decay are then combined, if the event includes no tracks other than those of the tag and the semileptonic candidate. Semileptonic decays are identified using the variable $`UE_{\mathrm{miss}}c\|\stackrel{}{p}_{\mathrm{miss}}\|`$, where $`E_{\mathrm{miss}}`$ and $`\stackrel{}{p}_{\mathrm{miss}}`$ are the missing energy and momentum of the $`D`$ meson decaying semileptonically. If the decay products of the semileptonic decay have been correctly identified, $`U`$ is expected to be zero, since only a neutrino is undetected. The resolution in $`U`$ is improved using constraints described in Ref. cleoc-neutral\_semilep . Due to the finite resolution of the detector, the distribution in $`U`$ is approximately Gaussian, centered at $`U=0`$ with $`\sigma 10\mathrm{MeV}`$. (The width varies by mode and is larger for modes with neutral pions.) To remove multiple candidates in each semileptonic mode, one combination is chosen per tag mode, based on the proximity of the invariant masses of the $`K_S^0`$, $`\overline{K}^0`$, $`\rho ^0`$, $`\pi ^0`$, or $`\omega `$ candidates to their expected masses. The yield for each semileptonic mode is determined from a fit to its $`U`$ distribution, as shown in Fig. 1 with all tag modes combined. In each case the signal is represented by a Gaussian and a Crystal Ball function cb to account for initial and final state radiation (FSR). The parameters describing the tails of the signal function are fixed with a GEANT-based Monte Carlo (MC) simulation GEANT . The background functions are determined by a MC simulation that incorporates all available data on $`D`$ meson decays. The backgrounds are small and arise mostly from misreconstructed semileptonic decays with correctly reconstructed tags rhoenu . The background shape parameters are fixed, while the background normalizations are allowed to float in all fits to the data. The mode $`D^+\omega e^+\nu _e`$ has never previously been observed. There are 8 events consistent with $`D^+\omega e^+\nu _e`$ in Fig. 1 (e). The background in the signal region ($`[25;+30]`$ MeV in $`U`$) is estimated to be $`0.4\pm 0.2`$ events. The probability for the background of 0.6 events to fluctuate to 8 or more events is $`2.4\times 10^7`$, which corresponds to significance exceeding $`5.0\sigma `$. Therefore, this is the first observation of $`D^+\omega e^+\nu _e`$. The absolute branching fractions in Table 2 are determined using $`=N_{\mathrm{signal}}/ϵN_{\mathrm{tag}}`$, where $`N_{\mathrm{signal}}`$ is the number of fully reconstructed $`D^+D^{}`$ events obtained by fitting the $`U`$ distribution, $`N_{\mathrm{tag}}`$ is the number of events with a reconstructed tag, and $`ϵ`$ is the effective efficiency for detecting the semileptonic decay in an event with an identified tag. A MC simulation where the relative population of tag yields across tag modes approximates that in the data is used to determine the efficiency. We have considered the following sources of systematic uncertainty and give our estimates of their magnitudes in parentheses. The uncertainties associated with the efficiency for finding a track (0.7%), $`\pi ^0`$ (2.0% for $`D^+\omega e^+\nu _e`$ and 4.3% for $`D^+\pi ^0e^+\nu _e`$) and $`K_S^0`$ (3.0%) are estimated using missing mass techniques with the data cleoc-Dtagging . Details on the uncertainties associated with positron identification efficiency (1.0%) are provided in Ref. cleoc-neutral\_semilep . Uncertainties in the charged pion and kaon identification efficiencies (0.3% per pion and 1.3% per kaon) are estimated using hadronic $`D`$ meson decays. The uncertainty in the number of tags (1.1%) is estimated by using alternative signal functions in the fits to the $`M_{\mathrm{bc}}`$ distributions and by varying the end point of the background function argus . The uncertainty in modeling the background shapes in the fits to the $`U`$ distributions (0.4% to 3.3% by mode) has contributions from the uncertainties in the simulation of the positron and hadron fake rates as well as the input branching fractions in the MC simulation. The uncertainty associated with the requirement that there be no additional tracks in tagged semileptonic events (0.3%) is estimated by comparing fully reconstructed $`D\overline{D}`$ events in data and MC. The uncertainty in the semileptonic reconstruction efficiencies due to imperfect knowledge of the semileptonic form factors is estimated by varying the form factors in the MC simulation within their uncertainties (1.0%) for all modes except $`D^+\rho ^0e^+\nu _e`$ and $`D^+\omega e^+\nu _e`$; for these a conservative uncertainty (3.0%) is taken, as no experimental information on the form factors in Cabibbo-suppressed pseudoscalar-to-vector transitions exists. The uncertainty associated with the simulation of FSR and bremsstrahlung radiation in the detector material (0.6%) is estimated by varying the amount of FSR modeled by the PHOTOS algorithm PHOTOS and by repeating the analysis with and without recovery of photons radiated by the positron. The uncertainty associated with the simulation of initial state radiation ($`e^+e^{}D\overline{D}\gamma `$) is negligible. There is a systematic uncertainty due to finite MC statistics (0.7% to 4.0% by mode). Non-resonant semileptonic decays $`D^+K^{}\pi ^+e^+\nu _e`$ are background to $`D^+\overline{K}^0e^+\nu _e`$. There is evidence from the FOCUS experiment for a non-resonant component consistent with an S-wave amplitude interfering with $`D^+\overline{K}^0e^+\nu _e`$ focusSWave . Its contribution, estimated to be 2.4% in this analysis, is subtracted in the calculation of the branching fraction of $`D^+\overline{K}^0e^+\nu _e`$ nonResSearch . Systematic uncertainties associated with the subtraction are due to imperfect knowledge of the amplitude and phase of the non-resonant component (1.0%), and its effect on the reconstruction efficiency (1.5%) swave\_effs . The lineshapes for semileptonic modes with wide resonances are simulated using a relativistic Breit-Wigner with a Blatt-Weisskopf form factor. A systematic uncertainty associated with the $`\overline{K}^0`$ lineshape (1.2%) is assigned by comparing the ($`K^{}\pi ^+`$) invariant mass distribution in the data to alternative lineshapes and the non-resonant contribution. For $`D^+\rho ^0e^+\nu _e`$, there is insufficient data to constrain the non-resonant background or the resonance lineshape. The systematic uncertainties from these two sources are expected to be much smaller than the current statistical uncertainty for this mode, and are neglected. These estimates of systematic uncertainty are added in quadrature to obtain the total systematic uncertainty (Table 2): 4.2%, 5.6%, 4.1%, 6.2%, and 7.8% for $`D^+\overline{K}^0e^+\nu _e`$, $`D^+\pi ^0e^+\nu _e`$, $`D^+\overline{K}^0e^+\nu _e`$, $`D^+\rho ^0e^+\nu _e`$, and $`D^+\omega e^+\nu _e`$, respectively. We now discuss the results presented in this Letter and the $`D^0`$ semileptonic study in Ref. cleoc-neutral\_semilep . The measured equality of the inclusive semileptonic widths of $`D^0`$ and $`D^+`$ mesons demonstrates that the source of the lifetime difference between them is attributable to differences in the hadronic widths. The widths of the isospin conjugate exclusive semileptonic decay modes of the $`D^0`$ and $`D^+`$ are related by isospin invariance of the hadronic current. The results obtained here and in Ref. cleoc-neutral\_semilep allow the most precise tests so far. The ratio $`\frac{\mathrm{\Gamma }(D^0K^{}e^+\nu _e)}{\mathrm{\Gamma }(D^+\overline{K}^0e^+\nu _e)}`$ is expected to be unity. The world average value is $`1.35\pm 0.19`$ PDG . Using our results and the lifetimes of the $`D^0`$ and $`D^+`$ PDG , we obtain: $`\frac{\mathrm{\Gamma }(D^0K^{}e^+\nu _e)}{\mathrm{\Gamma }(D^+\overline{K}^0e^+\nu _e)}=1.00\pm 0.05(\mathrm{stat})\pm 0.04(\mathrm{syst}).`$ The result is consistent with unity and with two recent less precise results: a measurement from BES II using the same technique BESII\_ratio and an indirect measurement from FOCUS FOCUS\_ratio ; FOCUS\_ksmunu . Ratios of isospin conjugate decay widths for other semileptonic decay modes are given in Table 3. As the data are consistent with isospin invariance, the precision of each branching fraction can be improved by averaging the $`D^0`$ and $`D^+`$ results for isospin conjugate pairs. The isospin-averaged semileptonic decay widths, with correlations among systematic uncertainties taken into account, are given in Table 4. The ratio of decay widths for $`D\pi e^+\nu `$ and $`DKe^+\nu `$ provides a test of the LQCD charm semileptonic rate ratio prediction unquenched\_LQCD . Using the isospin-averaged results in Table 4, we find $`\frac{\mathrm{\Gamma }(D^0\pi ^{}e^+\nu )}{\mathrm{\Gamma }(DKe^+\nu )}=(8.1\pm 0.7(\mathrm{stat})\pm 0.2(\mathrm{syst}))\times 10^2`$, consistent with LQCD and two recent results pikenu\_cleo ; pikenu\_focus . Furthermore, the ratio $`\frac{\mathrm{\Gamma }(DK^{}e^+\nu )}{\mathrm{\Gamma }(DKe^+\nu )}`$ is predicted to be in the range 0.5 to 1.1 (for a compilation see Ref. FOCUS\_ratio ). Using the isospin averages in Table 4, we find $`\frac{\mathrm{\Gamma }(DK^{}e^+\nu )}{\mathrm{\Gamma }(DKe^+\nu )}=0.63\pm 0.03(\mathrm{stat})\pm 0.02(\mathrm{syst})`$. Finally, summing all CLEO-c exclusive semileptonic branching fractions gives $`(D_{\mathrm{excl}}^0)=(6.1\pm 0.2(\mathrm{stat})\pm 0.2(\mathrm{syst}))`$% and $`(D_{\mathrm{excl}}^+)=(15.1\pm 0.5(\mathrm{stat})\pm 0.5(\mathrm{syst}))`$%. These are consistent with the world average inclusive semileptonic branching fractions: $`(D^0e^+X)=(6.9\pm 0.3)`$% and $`(D^+e^+X)=(17.2\pm 1.9)`$PDG , excluding the possibility of additional semileptonic modes of the $`D^0`$ and $`D^+`$ with large branching fractions. In summary, we have presented the most precise measurements to date of the absolute branching fractions of $`D^+`$ decays to $`\overline{K}^0e^+\nu _e`$, $`\pi ^0e^+\nu _e`$, $`\overline{K}^0e^+\nu _e`$, and $`\rho ^0e^+\nu _e`$, and the first observation and absolute branching fraction measurement of $`D^+\omega e^+\nu _e`$. We have combined these with measurements in Ref. cleoc-neutral\_semilep , which use the same data and analysis technique, to demonstrate that charm exclusive semileptonic decays are consistent with isospin invariance and to test other theoretical predictions. A comparison of the world average inclusive semileptonic branching fractions to the sum of the semileptonic branching fractions in this work excludes the possibility of additional semileptonic modes with large branching fractions. The precision achieved in this analysis is consistent with the expected performance of CLEO-c. CESR is currently running to collect a much larger $`\psi (3770)`$ data sample. It is expected that this sample will result in greatly improved measurements of $`D^0`$ and $`D^+`$ semileptonic branching fractions, measurements of the decay form factors, which are stringent tests of LQCD, and the CKM matrix elements $`V_{cs}`$ and $`V_{cd}`$ cleoc . We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. This work was supported by the National Science Foundation and the U.S. Department of Energy.
warning/0506/hep-th0506051.html	ar5iv	text	# 1 Introduction and General Ideas ## 1 Introduction and General Ideas During the sixties the study of hadronic particles was a mainstream area of theoretical Physics. The Regge trajectories were proposed and by the systematic study of scattering of Hadrons, the dualities between the $`s`$ and $`t`$ channels amplitudes were investigated and nicely realized by the Veneziano amplitude, The dual models, precursors of the modern String Theory, were one of the best candidates to explain the relevant Physics and provided the tools to explore the issues mentioned above. Despite some success, the experiments realized during the late sixties involving scatterings of large $`s`$ and large $`t`$ Mandelstam variables, but keeping fixed $`\frac{s}{t}`$ (fixed angle processes), gave results that showed that the amplitude was falling as a power law instead of the exponential law predicted by the dual models. This situation, together with the advent of QCD (that correctly predicted the scaling for fixed angle scattering and many other things), lead to the demise of the dual models for the study of hadronic Physics. Even when the original motivation to study dual models momentarily dissapeared, their rich structure kept many physicists interested and, after some technical subtleties were understood, this was the beginning of String Theory. As is known, after thirty years, string theorists have come back to the study of problems related to the hadronic world. Indeed, guided by the Maldacena Conjecture and their refinements , there have been very many interesting achievements in the area. They correspond to field theories with different amount of SUSY (including no SUSY) but “very similar” to QCD. The important point is that many characteristic features of QCD, like confinement, chiral symmetry breaking, etc; have been understood based on dual String theory backgrounds. In this paper, we study glueballs in some of the models mentioned above as being “very similar” to QCD (even when the models we will deal with here and those available in the literature, perhaps are not in the same universality class of QCD). Let us motivate a little bit the study of these glue-composed excitations. We know that the main distinction between a field theory in a confining phase and the same field theory in the Higgs phase is the presence of Regge trajectories, that do not occur in theories with Coulomb of Yukawa interactions. These Regge trajectories appear when plotting the spin $`J`$ and the squarred mass $`m^2`$ of the excitation, thus giving relation of the form $`J=\alpha ^{}m^2+\alpha _0,`$ with $`\alpha ^{}(1GeV)^2`$; this relation does not have in principle, an upper bound in $`J,m`$. It is due to these infinite number of Regge resonances, being interchanged in the $`s,t`$ channels of any hadron scattering that the beautiful structure of duality appeared in the models above mentioned. The glueballs should be some of these Regge excitations (making up a full trajectory if mixing with quarks is neglected) and this is a possible motivation to study them. From a modern QCD perspective, it is known that the cloud of gluons is what logically connects between a current quark (with mass of a few $`MeV`$) to a constituent quark, with mass of around 300 $`MeV`$. Since glue is part of the hadronic matter, we can consider color singlet composites of the form $`q\overline{q}g,gg,ggg`$ (apart from the mesons, baryons and exotics). The glueballs are composites made out of constituent glue, with no quark content. Of course, since we live in a world with quarks, one might think that the proposal of pure glue objects is impossible to study, because the quarks should run in loops when doing corrections to the operators, that leads to glueballs mixing with mesons, rendering the object not-pure glue. But lattice theorist (working in the quenched approximation) are not stopped by this. Indeed they took advantage of the limitation and have taught us many things about glueballs. Among the things that Lattice showed about QCD glueballs, we can mention the facts that: * there is a bound state spectrum * the lightest glueball is a scalar * the next is a tensor, 1.6 times heavier * the mass of the lightest glueball should be around 1630 $`MeV`$; see for example for a nice and clear review of these results. How are these lattice predictions experimentally checked? Experimentalist look for processes rich in glue production, like the $`J/\psi `$ decay, where the $`c,\overline{c}`$ quarks annihilate into gluons. Other process might be the $`p\overline{p}`$ annihilation, in this case the idea is that the quarks and anti-quarks in the initial hadrons annihilate completely, producing glue that later decays into hadrons. There are many glueballs candidates. One of them seems to be well established and is called $`f_0(1500)`$ with a width of $`112MeV`$ . From the string theory view point, using the Maldacena duality for the case of confining backgrounds, the study of glueballs (in the field theory dual to the background) proceeds by finding bound states for the fluctuations of the supergravity fields. Basically, the idea is to fluctuate all the fields in a given IIA, or IIB solution dual to a confining field theory, and linearizing in the fluctuated fields, study their eqs of motion (that are the Einstein, Maxwell and Bianchi eqs). The system is reduced to a Schroedinger problem. When solved, has eigenfunctions that we identify with the glueballs and eigenvalues that are identified with their masses. The fluctuations of the fields are dual to different operators in the gauge theory and what we are actually computing in the field theory side is the two-point correlation function of two glueball operators that should behave in a Wilson expansion as $$<O(x)O(y)>=\underset{j}{}c_je^{M_j\|xy\|}$$ where $`M_j`$ are the glueball masses. The quantum numbers of the glueballs $`J^{PC}`$ are determined on the basis of the spin ($`J`$) of the supergravity field and the R-symmetry quantum numbers (in a KK-harmonics decomposition) as studied, for example, in . We should point that this procedure is not totally clear in many of the available confining-models and it should be important to understand it better. This machinery has been applied to some confining models. Let us add that, since many of the existing Supergravity models are duals to confining field theories with only adjoint matter content, the objects under study are only glueballs (no hybrids) and since we work in the large $`N_c`$ regime, the glueballs are stable. Let us briefly review what was done in this subject. the original idea, described above, has been proposed by Witten in . Many papers followed, exploring this nice idea in different contexts. For example, confining models using black hole geometries were developed for $`QCD_3`$ and $`QCD_4`$ in , . Also, models based on rotating branes were introduced and other models based on $`AdS_5`$ with no SUSY . All these models have an spectrum that is numerically very close to the one obtained by Lattice methods. We should stress, that even when the comparisons between the lattice and “AdS” based results seem so accurate and promising, these calculations are done in opposite regimes. Indeed, the gravity-dual computation is in strong ’t Hooft coupling and this limitation is imposed by the Supergravity approximation. On the other hand, the Lattice computations are done at weak coupling, this seems to be a necessity of having a continuum limit because the lattice spacing $`a`$, has a relation $`a\mathrm{\Lambda }_{QCD}e^{1/g^2N}`$ with the QCD coupling and scale (other regularizations give similar results). One might think about doing strong coupling lattice computations, but they do not seem to be smoothly related to the continuum theory. It is possible that the numerical coincidences aluded above, are based on some dynamical principle to be understood. There exist a set of Supergravity models that preserve $`N=1`$ SUSY that have been object of lots of study and amusing advances. One of the models was put forward by Klebanov and Strassler and the glueballs in this model were carefully studied in the set of papers ,. The results indicate that, for the Klebanov-Strassler model, the masses of the $`0^{++},\mathrm{\hspace{0.33em}1}^{},\mathrm{\hspace{0.33em}2}^{++}`$ in the strong ’t Hooft coupling limit, fall in a linear trajectory. Also, the finding of a massless excitation showed that this cascading field theory is not in the same universality class of QCD, because of the reasons we explained above. <sup>1</sup><sup>1</sup>1We thank Oliver Jahn for extensive discussions on many of the points touched in this introduction ### 1.1 Motivations and organization of this paper We mentioned above a set of Supergravity duals to confining models and up to this point we just commented on the one proposed by Klebanov and Strassler . There exist some other models that are based on wrapped D-branes. The main idea here is to consider the low energy field theory in $`(k+1)`$ dimensions, obtained by wrapping a Dp brane on a $`(pk)`$ cycle. Some subtleties of this type of models will be explained in section 2 of this paper. Here we just want to emphasize that the glueballs spectrum in this case is poorly understood. Indeed, there is a paper , where a study was initiated. We believe that this study is not completely correct from a technical viewpoint (we believe that incorrect eqs were used) and the conclusions expressed there are, even when intuitively understandable, also not totally correct. The main point of that paper is that in one of these models it is necessary to introduce a hard cut-off in order to have a discrete spectrum. In this paper we re-analize this statement and propose a different result, basically that all these models do have discrete spectrum of glueballs and give a way of computing it. The spectrum even though discrete is not normalizable. In order to get normalizable states we need to introduce a regularization procedure. The regularization we propose is not the introduction of a hard cut-off but is more in the spirit of the Wilson loop calculations , where a non-physical part is subtracted. In this paper we will not make much emphasis on the numerical aspects of the problem. Indeed, even when discrete states are numerically obtained, we will not worry here about comparisons with lattice results, that as explained above are perhaps not very significative. The main objective of this work is to study qualitative features of the spectrum, point out differences with previously studied cases, and propose a procedure of computing and regularizing in these wrapped branes set-ups. This paper is divided in two parts, one dealing with a particular type IIB model and the other with a type IIA model. Both parts have been written and can be read in parallel and almost independently. In section 2, we describe in detail the two models we will be using, one based in type IIB, with D5 branes wrapping a two-cycle inside the resolved conifold. The other in type IIA, based on D6 branes wrapping a three cycle in the deformed conifold. Section 3 deals with the glueballs in the type IIB set-up, while section 4 sketches the results corresponding to the type IIA model. We did not carry the problem to an end because of the need of a more precise numerical analysis, since the solution is only numerically known. Section 5 presents conclusions and possible future work proposed to the interested reader. There are very detailed Appendixes that carefully explain all the computations in sections 3 and 4. ## 2 $`𝒩=1`$ SYM models from wrapped branes In this section, we write an account of duals to N=1 SYM from wrapped branes. The two main models on which we will concentrate are the ones based on D5 branes wrapping a two cycle, that we will consider to be a two-sphere inside a CY3 fold and D6 branes on a three cycle (a squashed three sphere), also inside a CY3 fold. We will present the solutions in detail, and explain the main characteristics of the dual gauge theory. We will emphasize the existence of ‘extra’ modes called KK modes with mass of the same order of the confinement scale. Since our interest in this paper is on glueballs, we will discuss the influence of this ‘extra’ modes in the computation of glueballs for $`N=1`$ SYM. ### 2.1 D5 branes wrapping $`S^2`$ We will work with the model presented in (the solution was first found in a 4d context in ) and described and studied in more detail in the paper . Let us briefly describe the main points of this supergravity dual to $`N=1`$ SYM and its UV completion. Suppose that we start with N $`D5`$ branes, the field theory living on them is $`(5+1)`$SYM with 16 supercharges. Then, suppose that we wrap two directions of the D5 branes on a curved two manifold that can be choosen to be a sphere. In order to preserve SUSY a twisting procedure has to be implemented. The one we will be interested in this section, deals with a twisting that preserves four supercharges. In this case the two-cycle mentioned above lives inside a CY3 fold. Notice that this supergravity solution will be dual to a four dimensional field theory, only for low energies (small values of the radial coordinate). Indeed, at high energies, the modes of the gauge theory start to explore the two cycle and the theory becomes first N=1 SYM in six dimensions and then, the blowing-up of the dilaton forces us to S-dualize and a little string theory completes the model in the UV. In this sense, to study only the 4d-SYM part of the background, a procedure that “substracts” the unwanted UV completion, should be useful. We will elaborate on this in Section 3.1. The supergravity solution corresponding to the case of interest in this section, the one preserving four supercharges, has the topology of $`R^{1,3}\times R\times S^2\times S^3`$ and there is a fibration between the two spheres that allows the SUSY preservation. The topology of the metric, near $`r=0`$ is $`R^{1,6}\times S^3`$. The full solution and Killing spinors are written in detail in . Let us revise it here for reference. The metric in Einstein frame reads, $$ds_{10}^2=\alpha ^{}g_sNe^{\frac{\varphi }{2}}\left[\frac{1}{\alpha ^{}g_sN}dx_{1,3}^2+e^{2h}\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right)+dr^2+\frac{1}{4}(w^iA^i)^2\right],$$ (2.1) where $`\varphi `$ is the dilaton. The angles $`\theta [0,\pi ]`$ and $`\phi [0,2\pi )`$ parametrize a two-sphere. This sphere is fibered in the ten dimensional metric by the one-forms $`A^i`$ $`(i=1,2,3)`$. Their expression can be written in terms of a function $`a(r)`$ and the angles $`(\theta ,\phi )`$ as follows: $$A^1=a(r)d\theta ,A^2=a(r)\mathrm{sin}\theta d\phi ,A^3=\mathrm{cos}\theta d\phi .$$ (2.2) The $`w^i`$’s appearing in eq. (2.1) are the $`su(2)`$ left-invariant one-forms, satisfying $`w^1`$ $`=`$ $`\mathrm{cos}\psi d\stackrel{~}{\theta }+\mathrm{sin}\psi \mathrm{sin}\stackrel{~}{\theta }d\stackrel{~}{\phi },`$ $`w^2`$ $`=`$ $`\mathrm{sin}\psi d\stackrel{~}{\theta }+\mathrm{cos}\psi \mathrm{sin}\stackrel{~}{\theta }d\stackrel{~}{\phi },`$ $`w^3`$ $`=`$ $`d\psi +\mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\phi }.`$ (2.3) The three angles $`\stackrel{~}{\phi }`$, $`\stackrel{~}{\theta }`$ and $`\psi `$ take values in the rank $`0\stackrel{~}{\phi }<2\pi `$, $`0\stackrel{~}{\theta }\pi `$ and $`0\psi <4\pi `$. For a metric ansatz such as the one written in (2.1) one obtains a supersymmetric solution when the functions $`a(r)`$, $`h(r)`$ and the dilaton $`\varphi `$ are: $`a(r)`$ $`=`$ $`{\displaystyle \frac{2r}{\mathrm{sinh}2r}},`$ $`e^{2h}`$ $`=`$ $`r\mathrm{coth}2r{\displaystyle \frac{r^2}{\mathrm{sinh}^22r}}{\displaystyle \frac{1}{4}},`$ $`e^{2\varphi }`$ $`=`$ $`e^{2\varphi _0}{\displaystyle \frac{2e^h}{\mathrm{sinh}2r}},`$ (2.4) where $`\varphi _0`$ is the value of the dilaton at $`r=0`$. Near the origin $`r=0`$ the function $`e^{2h}`$ behaves as $`e^{2h}r^2`$ and the metric is non-singular. The solution of the type IIB supergravity includes a Ramond-Ramond three-form $`F_{(3)}`$ given by $$\frac{\alpha ^{}}{N}F_{(3)}=\frac{1}{4}\left(w^1A^1\right)\left(w^2A^2\right)\left(w^3A^3\right)+\frac{1}{4}\underset{a}{}F^a\left(w^aA^a\right),$$ (2.5) where $`F^a`$ is the field strength of the su(2) gauge field $`A^a`$, defined as: $$F^a=dA^a+\frac{1}{2}ϵ_{abc}A^bA^c.$$ (2.6) The different components of $`F^a`$ are: $$F^1=a^{}drd\theta ,F^2=a^{}\mathrm{sin}\theta drd\phi ,F^3=(\mathrm{\hspace{0.17em}1}a^2)\mathrm{sin}\theta d\theta d\phi ,$$ (2.7) where the prime denotes derivative with respect to $`r`$. Since $`dF_{(3)}=0`$, one can represent $`F_{(3)}`$ in terms of a two-form potential $`C_{(2)}`$ as $`F_{(3)}=dC_{(2)}`$. Actually, it is not difficult to verify that $`C_{(2)}`$ can be taken as: $`{\displaystyle \frac{\alpha ^{}C_{(2)}}{N}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}[\psi (\mathrm{sin}\theta d\theta d\phi \mathrm{sin}\stackrel{~}{\theta }d\stackrel{~}{\theta }d\stackrel{~}{\phi })\mathrm{cos}\theta \mathrm{cos}\stackrel{~}{\theta }d\phi d\stackrel{~}{\phi }`$ (2.8) $`a(d\theta w^1\mathrm{sin}\theta d\phi w^2)].`$ Moreover, the equation of motion of $`F_{(3)}`$ in the Einstein frame is $`d\left(e^\varphi {}_{}{}^{}F_{(3)}^{}\right)=0`$, where $``$ denotes Hodge duality. Let us stress here that the previous configuration is non-singular. Finally, let us comment on the fact that the BPS equations also admit a solution in which the function $`a(r)`$ vanishes, i.e. in which the one-form $`A^i`$ has only one non-vanishing component, namely $`A^3`$. We will refer to this solution as the “abelian” (or “singular”) $`𝒩=1`$ background. Its explicit form can be easily obtained by taking the $`r\mathrm{}`$ limit of the functions given in eq. (2.4). Notice that, indeed $`a(r)0`$ as $`r\mathrm{}`$ in eq. (2.4). Neglecting exponentially suppressed terms, one gets: $$e^{2h}=r\frac{1}{4},(a=0),$$ (2.9) while $`\varphi `$ can be obtained from the last equation in (2.4). The metric of the abelian background is singular at $`r=1/4`$ (the position of the singularity can be moved to $`r=0`$ by a redefinition of the radial coordinate). This IR singularity of the abelian background is removed in the non-abelian metric by switching on the $`A^1,A^2`$ components of the one-form (2.2). ### 2.2 Some analysis of this model Let us first summarize the field theory aspects of the dual to the gravity solution we will be mainly concerned with. The main characteristic is that it contains a four dimensional Minkowski space, a radial direction and a two sphere fibered over a three sphere. In this solution was argued to be dual to $`N=1`$ SYM. Let us analize the claim a little more, the field theory at low energies (low compared to the inverse size of the two-sphere) has degrees of freedom given by a vector field and a Majorana spinor (in 4d). When increasing in energy, other modes with mass of the order of the inverse size of the $`S^2`$ appear in the spectrum. These are called KK modes and can be seen as coming from the reduction of the maximally D5 branes SUSY field theory on a two dimensional sphere and a twisting (explained below) are performed. When the energy is high enough the excitations of the theory propagate in $`5+1`$ dimensions and the UV-completion of our minimally SUSY four dimensional field theory is the six dimensional little string theory living on $`N`$ NS5 branes. Let us recall briefly the twisting procedure. One has a $`5+1`$ field theory (that lives on $`N`$ D5 branes) that has gauge fields, fermions and four scalars, all in the adjoint of the $`SU(N)`$ gauge group. We rewrite the $`SO(1,5)\times SU(2)_L\times SU(2)_R`$ group quantum numbers of the fields above in terms of $`SO(1,3)\times SO(2)\times SU(2)_L\times SU(2)_R`$ and then we mix the quantum numbers respect to $`SO(2)`$ with those of another $`SO(2)`$ that lives inside one of the $`SU(2)^{}s`$. After this twisting procedure is performed, we are left with fields that under $`SO(1,3)\times U(1)\times SU(2)`$ transform as $$A_\mu ^a=(4,0,1),\mathrm{\Phi }^a=(1,\pm ,1),\varphi ^a=2(1,\pm ,2).$$ (2.10) for the bosons that can be seen to be a massless gauge field, a massive scalar (coming from the gauge field) and other massive scalars (that originally represented the positions of the D5 branes in $`R^4`$). As a general rule, all the fields that do transform under the twisted $`U(1)`$, the second entry in the charges above, will be massive. For the fermions we will have $$\psi ^a=(2,0,1),(\overline{2},0,1),(2,++,1),(\overline{2},,1),(2,0,2),(\overline{2},0,2)$$ (2.11) that is a Majorana spinor in four dimensions that is massless and then we have massive ones (those whose quantum number under the twisted $`U(1)`$ is not zero). The KK modes are the massive modes mentioned above. Their mass is of the order $`M_{KK}^2=(R_{S2})^2=\frac{1}{g_s\alpha ^{}N}`$. The dynamics of these KK modes, mixes with the dynamics of confinement in this model, because the scale of strong coupling of the theory is of the order of the KK mass. If we could work with a sigma model for the string in this background (or in the S-dual NS5 background) to all orders in $`\alpha ^{}`$ we could decouple both behaviours. Meanwhile, the dynamics of these KK modes has not been studied in great detail, but some progress have been made, for example in the papers . Finally, we would like to mention a paper where a very careful study of the KK modes spectrum have been done, also pointing a coincidence with $`N=1^{}`$ theory in a given Higgs vacuum . Let us briefly comment on the influence of these KK modes on the glueballs spectrum. Indeed, once the strong coupling regime of the field theory is attained, one possible way to compute is using these supergravity backgrounds. Given that we are in the supergravity approximation, the spectrum of our model includes these KK ‘contaminations’ (this feature repeats in all the dual to non-conformal field theories). Obviously, our glueballs will be of two types, those coming from condensates of the gluon and gluino, those ‘composed’ out of KK modes and finally, hybrids, composed out of SYM fields and KK modes. We would like to discard those with some KK constituent. We will comment in the conclusion section on a possibility to do this. Finally, let us mention that there are many succesful checks showing that the supergravity background presented above captures different non-perturbative aspects of $`N=1`$ SYM. We will not discuss these many checks here, instead, we refer the interested reader to the very careful reviews . ### 2.3 D6 branes wrapping $`S^3`$ Now, let us comment on the models based on D6 branes wrapping a calibrated three-cycle inside a CY3 fold. The progress in this direction originated from the duality between Chern-Simons gauge theory on $`S^3`$ at large $`N`$ and topological string theory on a blown up Calabi-Yau conifold . This duality was embedded in string theory as a duality between the IIA string theory of $`N`$ D6-branes wrapping the blown up $`S^3`$ of the deformed conifold and IIA string theory on the small resolution of the conifold with $`N`$ units of two form Ramond-Ramond flux through the blown up $`S^2`$ and no branes . The D6-brane side of the duality involves an $`𝒩=1`$ gauge theory in four dimensions that is living on the non-compact directions of the branes, at energies that do not probe the wrapped $`S^3`$. Just like before, in order for the wrapped branes to preserve some supersymmetry, one has to embedd the spin connection of the wrapped cycle into the gauge connection, which is known as twisting the theory. When we have flat D6 branes, the symmetry group of the configuration is $`SO(1,6)\times SO(3)_R`$. The spinors transform in the (8,2) of the isometry group and the scalars in the (1,3), whilst the gauge particles are in the (7,1) . Wrapping the D6 brane on the three-sphere breaks the group to $`SO(1,3)\times SO(3)\times SO(3)_R`$. The technical meaning of twisting is that the two $`SO(3)`$s get mixed to allow the existence of four dimensional spinors that transform as scalars under the new twisted $`SO(3)`$ . One can then see that the remaining particles in the spectrum that transform as scalars under the twisted $`SO(3)`$ are the gauge field and four of the initial sixteen spinors. Thus the massless field content is that of $`𝒩=1`$ SYM. Like in the model analyzed in the previous section, apart from these fields, there will be massive modes, whose mass scale is set by the size of the curved cycle. When we probe the system with very low energies, we find only the spectrum of $`𝒩=1`$ SYM. For D6 branes in flat space, the ‘decoupling’ limit does not completely decouple the gauge theory modes from bulk modes . In our case, we expect a good gauge theory description only when the size of the wrapped three-cycle is large, which implies that we have to probe the system with very low energies to get 3+1 dimensional SYM . In this case, the size of the two cycle in the flopped geometry is very near to zero, so a good gravity description is not expected. In short, we must keep in mind that the field theory we will be dealing with has more degrees of freedom than pure $`𝒩=1`$ SYM, thus the glueballs masses that one might obtain following the procedure explained in the following sections might be ‘contamined’ by glueballs composed out of KK modes or hybrids composed out of KK modes and gluons or gauginos. Again, how to decouple the ones we are interested into from those glueballs ‘composed’ of KK modes is going to be discussed in the conclusion section. Finally, let us add that the duality described above is naturally understood by considering M-theory on a $`G_2`$ holonomy metric . In eleven dimensions, $`G_2`$ holonomy implements $`𝒩=1`$ as pure gravity. One starts with a singular $`G_2`$ manifold that on dimensional reduction to IIA string theory corresponds to $`N`$ D6 branes wrapping the $`S^3`$ of the deformed conifold. There is an $`SU(N)`$ gauge theory at the singular locus/D6 brane. This configuration describes the UV of the gauge theory. As the coupling runs to the IR, a blown up $`S^3`$ in the $`G_2`$ manifold shrinks and another has fixed size. This flop is smooth in M-theory physics. The metrics will be discussed in more detail in the following sections. In the IR regime, the $`G_2`$ manifold is non-singular and dimensional reduction to IIA gives precisely the aforementioned small resolution of the conifold with no branes and RR flux. Let us now, write explicitly the background on which we will be interested. It is conveninet to start with the eleven dimensional M-theory background, that reads $$ds_{11}^2=dx_{1,3}^2+ds_7^2$$ (2.12) with $`ds_7^2=dr^2+a(r)^2\left[(\mathrm{\Sigma }_1+g(r)\sigma _1)^2+(\mathrm{\Sigma }_2+g(r)\sigma _2)^2\right]+c(r)^2(\mathrm{\Sigma }_3+g_3(r)\sigma _3)^2`$ $`+b(r)^2\left[\sigma _1^2+\sigma _2^2\right]+f(r)^2\sigma _3^2,`$ (2.13) where $`\mathrm{\Sigma }_i,\sigma _i`$ are left-invariant one-forms on the $`SU(2)`$s (2.3). The six functions are not all independent $$g(r)=\frac{a(r)f(r)}{2b(r)c(r)},g_3(r)=1+2g(r)^2.$$ (2.14) None of the radial functions are known explicitly, although the asymptotics at the origin and at infinity are known. The asymptotics are found by finding Taylor series solutions to the first order equations for the radial functions. The equations are , $`\dot{a}={\displaystyle \frac{c}{2a}}+{\displaystyle \frac{a^5f^2}{8b^4c^3}},`$ $`\dot{b}={\displaystyle \frac{c}{2b}}{\displaystyle \frac{a^2(a^23c^2)f^2}{8b^3c^3}},`$ $`\dot{c}=1+{\displaystyle \frac{c^2}{2a^2}}+{\displaystyle \frac{c^2}{2b^2}}{\displaystyle \frac{3a^2f^2}{8b^4}},`$ $`\dot{f}={\displaystyle \frac{a^4f^3}{4b^4c^3}}.`$ (2.15) As $`r0`$ one has $`a(r)`$ $`=`$ $`{\displaystyle \frac{r}{2}}{\displaystyle \frac{(q_0^2+2)r^3}{288R_0^2}}{\displaystyle \frac{(7429q_0^2+31q_0^4)r^5}{69120R_0^4}}+\mathrm{},`$ $`b(r)`$ $`=`$ $`R_0{\displaystyle \frac{(q_0^22)r^2}{16R_0}}{\displaystyle \frac{(1321q_0^2+11q_0^4)r^4}{1152R_0^3}}+\mathrm{},`$ $`c(r)`$ $`=`$ $`{\displaystyle \frac{r}{2}}{\displaystyle \frac{(5q_0^28)r^3}{288R_0^2}}{\displaystyle \frac{(232353q_0^2+157q_0^4)r^5}{34560R_0^4}}+\mathrm{},`$ $`f(r)`$ $`=`$ $`q_0R_0+{\displaystyle \frac{q_0^3r^2}{16R_0}}+{\displaystyle \frac{q_0^3(14+11q_0^2)r^4}{1152R_0^3}}+\mathrm{},`$ (2.16) where $`q_0`$ and $`R_0`$ are constants. Note that $`a(r)`$ and $`c(r)`$ collapse and the other two functions do not. As $`r\mathrm{}`$ we have $`a(r)`$ $`=`$ $`{\displaystyle \frac{r}{\sqrt{6}}}{\displaystyle \frac{\sqrt{3}q_1R_1}{\sqrt{2}}}+{\displaystyle \frac{(27\sqrt{6}96h_1)R_1^2}{96r}}+\mathrm{},`$ $`b(r)`$ $`=`$ $`{\displaystyle \frac{r}{\sqrt{6}}}{\displaystyle \frac{\sqrt{3}q_1R_1}{\sqrt{2}}}+{\displaystyle \frac{h_1R_1^2}{r}}+\mathrm{},`$ $`c(r)`$ $`=`$ $`{\displaystyle \frac{r}{3}}+q_1R_1{\displaystyle \frac{9R_1^2}{8r}}+\mathrm{},`$ $`f(r)`$ $`=`$ $`R_1{\displaystyle \frac{27R_1^3}{8r^2}}{\displaystyle \frac{81R_1^4q_1}{4r^3}}+\mathrm{}.`$ (2.17) With constants $`R_1,q_1,h_1`$. Note that $`f(r)`$ stabilises. Three constants appear to this order, whilst there were only two constants in the expansion around the origin. This just means that for some values of these constants, the corresponding solution will diverge before it reaches zero. In any case, we find no $`h_1`$ dependence in the results below. We can reduce this to Type IIa and we will find a non-singular background with dilaton, metric and RR one form excited, that reads, $`ds_{IIA,string}^2=2e^{2/3\varphi }(dx_{1,3}^2+dr^2+b(r)^2(\sigma _1^2+\sigma _2^2)+a(r)^2((\mathrm{\Sigma }_1+g(r)\sigma _1)^2+(\mathrm{\Sigma }_2+g(r)\sigma _2)^2)+`$ $`{\displaystyle \frac{f^2c^2}{f^2+c^2(1+g_3)^2}}(\mathrm{\Sigma }_3\sigma _3)^2)`$ $`4e^{4/3\varphi }=f(r)^2+c(r)^2(1+g_3(r))^2,A_1=\mathrm{cos}\theta d\phi +\mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\phi }+{\displaystyle \frac{f^2c^2(1g_3^2)}{f^2+c^2(1+g_3)^2}}(\sigma _3\mathrm{\Sigma }_3)`$ (2.18) Where we have defined $`dx_{11}=d\psi ^{}+d\psi `$ and $`d\widehat{\psi }=d\psi d\psi ^{}`$. So, to summarize the things clearly, let us write the metric of our IIA solution in Einstein frame (as will be used below), $`ds_E^2=2e^{\varphi /6}(dx_{1,3}^2+dr^2+(b(r)^2+g(r)^2)(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)+a(r)^2(d\stackrel{~}{\theta }^2+\mathrm{sin}^2\stackrel{~}{\theta }d\stackrel{~}{\phi }^2)+`$ $`+2g(r)a(r)^2[\mathrm{cos}\widehat{\psi }(d\theta d\stackrel{~}{\theta }+\mathrm{sin}\theta \mathrm{sin}\stackrel{~}{\theta }d\phi d\stackrel{~}{\phi })+\mathrm{sin}\widehat{\psi }(\mathrm{sin}\theta d\stackrel{~}{\theta }d\phi sin\stackrel{~}{\theta }d\theta d\stackrel{~}{\phi })]+`$ $`+{\displaystyle \frac{f^2c^2}{f^2+c^2(1+g_3)^2}}(d\widehat{\psi }+\mathrm{cos}\theta d\phi \mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\phi })^2)`$ (2.19) with the same dilaton as in (2.18), besides, the field strength $`F_2`$ reads, $$F_2=k^{}(r)dr(d\widehat{\psi }+\mathrm{cos}\theta d\phi \mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\phi })(k(r)+1)\mathrm{sin}\theta d\theta d\phi +(k(r)1)\mathrm{sin}\stackrel{~}{\theta }d\stackrel{~}{\theta }d\stackrel{~}{\phi },$$ (2.20) and $$k(r)=\frac{f^2c^2(1g_3^2)}{f^2+c^2(1+g_3)^2}.$$ To end this section, let us briefly revise what checks exist of the duality between the backgrounds presented here and $`N=1`$ SYM. Of course, the number of supercharges match, there is a nice picture of confinement in terms of a Wilson loop computation in IIA. But most of the presently known matchings of $`N=1`$ SYM with $`G_2`$ holonomy M-theory come from considering membrane instantons as gauge theory instantons that generate the superpotential , membranes wrapped on one-cycles in the IR geometry that are super QCD strings in the gauge theory , and fivebranes wrapped on three-cycles that give domain walls in the gauge theory . These matchings above, are essentially topological and do not use the explicit form of the $`G_2`$ metrics. In the category of test/checks that use the form of the metric, we can mention , where rotating membranes in these $`G_2`$ backgrounds have been studied and relations for large operators in SYM have been reproduced. We should also mention , where a very nice picture of the chiral anomaly of SYM have been developed. Perhaps less promising is the fact that confining string tensions do not arise as cleanly in the IIA backgrounds as in the type IIB case . There are many aspects of this duality that are not on a very firm basis and we think that through study, these unclear points might become clear. The results that we will present in Section 4, use the explicit form of the metric and should be considered to belong to this second category of tests. ## 3 Glueballs from type IIB solution As we explained above, in order to study glueballs, we need the variation of the eqs of motion. In the type IIB case, for the solution of D5 branes wrapping $`S^2`$ inside a $`CY_3`$-fold discussed in the previous section, the Einstein eqs for the metric, dilaton and three form read, <sup>2</sup><sup>2</sup>2We will denote the contraction of indexes with $``$ symbols, so, for example $`g^{\mu \nu }A_{\mu kl}B_{\nu jp}=A_{\mu kl}B_{jp}^\mu =A_{kl}B_{jp}^{}`$. $$R_{\mu \nu }=\frac{1}{2}_\mu \varphi _\nu \varphi +\frac{g_se^\varphi }{4}\left(F_\mu F_\nu ^{}\frac{1}{12}g_{\mu \nu }F_3^2\right)$$ (3.21) by contracting we get the Ricci scalar eq, $$R=\frac{1}{2}(_\mu \varphi )^2+\frac{g_se^\varphi }{24}F_3^2.$$ (3.22) The dilaton, Maxwell and Bianchi eqs. are, $$^2\varphi =\frac{g_se^\varphi }{12}F_3^2,_\mu \left(\sqrt{g}e^\varphi F^{\mu \nu \rho }\right)=0,_{[\mu }F_{\kappa \nu \rho ]}=0$$ (3.23) Now, let us study the fluctuations of these eqs. Let us assume that the background fields vary according to $$g_{\mu \nu }g_{\mu \nu }+ϵh_{\mu \nu },\varphi \varphi +ϵ\delta \varphi ,F_{\mu \nu \rho }F_{\mu \nu \rho }+ϵ\delta F_{\mu \nu \rho }$$ (3.24) Keeping only linear order in the parameter $`ϵ`$, we get eqs for the fluctuated fields (that can be found written in detail in Appendix B). The metric fluctuation, $`h_{\mu \nu }`$ can be splitted in its worldvolume, internal and mixed parts: $`h_{ij},h_{\alpha \beta }`$ and $`h_{i\alpha }`$ respectively. We assume that there is no fluctuation in the mixed part, $`h_{i\alpha }=0`$ and perform a Weyl shift in the worldvolume fluctuation, <sup>3</sup><sup>3</sup>3 The standard Weyl shift in a D dimensional spacetime is $`\lambda /5=1/(D2)`$; this value is required to simplify the variation of $`^2`$. Here we choose to leave $`\lambda `$ as a constant to be determined later. $`h_{ij}=h_{ij}^{}+{\displaystyle \frac{\lambda }{5}}g_{ij}h_\alpha ^\alpha ,,h_{\alpha \beta }=h_{(\alpha \beta )}+{\displaystyle \frac{1}{5}}g_{\alpha \beta }h_a^a`$ (3.25) We denote with latin indices the worldvolume and transverse coordinates $`i,j=\stackrel{}{x},t,r`$ and with with greek indices the internal coordinates, $`\alpha ,\beta =\theta ,\stackrel{~}{\theta },\phi ,\stackrel{~}{\phi },\psi `$. Also, $`\lambda `$ is a constant and $`h_{(\alpha \beta )}`$ is the symmetric traceless part of the fluctuation in the internal directions. Notice that the trace of the metric fluctuation over all the space is $`h_\mu ^\mu =(\lambda +1)h_\alpha ^\alpha `$. From here on we will denote with $`x`$ the radial and ‘gauge theory’ coordinates and with $`y`$ the internal coordinates (the angles on the spheres). Imposing a de Donder and Lorentz type condition, $`^\alpha h_{(\alpha \beta )}=^\alpha h_{\alpha i}=0`$, the decomposition in harmonics for the fluctuated fields is, $`h_\alpha ^\alpha (x,y)=\mathrm{\Sigma }_I\mathrm{\Pi }^I(x)Y^I(y),h_{(\alpha \beta )}=\mathrm{\Sigma }_Ib^I(x)Y_{(\alpha \beta )}^I(y),`$ $`h_{ij}=\mathrm{\Sigma }_IH_{ij}(x)Y^I(y),\delta \varphi (x,y)=\mathrm{\Sigma }_If^I(x)Y^I(y)`$ $`C_{ij}=\mathrm{\Sigma }_Ia_{ij}^I(x)Y^I(y),C_{i\alpha }=\mathrm{\Sigma }_Ia_i^I(x)Y_\alpha ^I(y),C_{\alpha \beta }=\mathrm{\Sigma }_Ia^I(x)Y_{\alpha \beta }^I(x).`$ (3.26) Where $`C_2`$ is the two form potential. We set $`h_{ij}^{}=0`$. Given this ansatz the equations of motion can be consistently solved for the other fluctuations. Using the decomposition in harmonics, keeping only the s-wave and choosing a particular value for $`\lambda =\frac{35}{29}`$ the two equations for the fluctuations nicely combine into just one equation (again, see the Appendix B for details). This final eq. reads, $$^2f(x)+2g^{rr}_rf(x)_r\varphi +\frac{g_s}{12}e^\varphi \left[\frac{5}{2}F_3^2+\frac{293}{10}(g^{ab}F_aF_b^{}\frac{35}{29}g^{rr}F_rF_r^{})\right]f(x)=0$$ (3.27) Expanding (3.27) in plane waves, $`f(x)=F(r)e^{(iKx)}`$, we have, $`g^{rr}{\displaystyle \frac{d^2F(r)}{dr^2}}`$ $`+`$ $`\left(2g^{rr}_r\varphi +{\displaystyle \frac{1}{\sqrt{g}}}_r(\sqrt{g}g^{rr})\right){\displaystyle \frac{dF(r)}{dr}}`$ $`+`$ $`\left({\displaystyle \frac{g_s}{12}}e^\varphi [{\displaystyle \frac{5}{2}}F_3^2+{\displaystyle \frac{293}{10}}(g^{ab}F_aF_b^{}{\displaystyle \frac{35}{29}}g^{rr}F_rF_r^{})]K^2g^{xx}\right)F(r)=0`$ This equation will be solved numerically; the glueball masses are given by the eigenvalues $`K^2`$ for which there is a solution with appropriate boundary conditions. Before studying the boundary conditions of this problem let us cast eq. (LABEL:finalplane) in a more familiar way. To save us some writing, denote the coefficient of the first derivative term ($`\frac{dF}{dr}`$) in (LABEL:finalplane) as $`\mu (r)`$ and the coefficient of $`F(r)`$ as $`\alpha (r)`$, that is, $`\mu (r)=2_r\varphi +{\displaystyle \frac{g_{rr}}{\sqrt{g}}}_r(\sqrt{g}g^{rr})`$ $`\alpha (r)={\displaystyle \frac{g_sg_{rr}}{12}}e^\varphi [{\displaystyle \frac{5}{2}}F_3^2+{\displaystyle \frac{293}{10}}(g^{ab}F_aF_b^{}{\displaystyle \frac{35}{29}}g^{rr}F_rF_r^{})]`$ (3.29) Making a change of variables, $`(r)=e^{\frac{1}{2}^r\mu (z)𝑑z}F(r),`$ equation (LABEL:finalplane) can be written in a Schrödinger form $$\frac{d^2(r)}{dr^2}VS(r)=0,$$ (3.30) where $$VS(r)=\alpha (r)+\frac{1}{4}(\mu ^2(r)+2\frac{d\mu (r)}{dr})+K^2$$ (3.31) A graph of the Schrödinger potential VS(r) is given in Figure 1. At this point it is convenient to recall that, as explained in section two, the full model, considering its UV completion is not dual to a non-abelian gauge theory, but to a little string theory. Indeed, in the UV, due to the divergent dilaton, the solution has to be S-dualized yielding a IIB solution with NS five branes. In the decoupling limit this background is not dual to $`𝒩=1`$ SYM in four dimensions but to a higher dimensional, 5+1, little string theory. Naturally, trying to calculate observables by simply using the solution up to infinity might not yield sensible answers. In what follows we want to study the glueball spectrum and, if necessary, propose a regularization procedure. To calculate the $`0^{++}`$ mass we have to numerically find eigenvalues satisfying equation and appropriate boundary conditions. The asymptotic behavior of the potential is, $$VS(r\mathrm{})=\frac{92}{5}+K^2$$ (3.32) $$VS(r0)=\frac{1036}{45}+K^2$$ (3.33) Choosing the exponentially decreasing solution at infinity we get, $$(r\mathrm{})e^{\sqrt{\frac{92}{5}+K^2}r}.$$ (3.34) At the origin we demand a smooth solution, $$\frac{dF(r)}{dr}=0.$$ (3.35) Similarly to Klebanov-Strassler, satisfying the boundary conditions implies that the eigenvalues are bounded from below, $`K^2>92/5`$ and thus, there is a mass gap. But here, in addition of being bounded from below, the eigenvalues are also bounded from above, $`K^2<1036/45`$. A similar phenomenon was observed in in the context of non-commutative gauge theories. Another difference with the Klebanov-Strassler model is that here the boundary condition at infinity depends on the eigenvalue. In a technical sense, each eigenvalue defines a different problem -different boundary condition- and the spectrum is then given by the eigenvalues of this collection of problems; It is a more general situation than the standard eigenvalue problem. Using the WKB method we can estimate the eigenvalues $`K^2`$ for which there exists a solution of (3.30) satisfying (3.34) and(3.35). Also, it can be shown numerically that the WKB integral $`_0^{r(K^2)}\sqrt{VS(r)}𝑑r`$ is a monotonically decreasing function of $`K^2`$. And this fact can be used to prove that there is only one eigenvalue in the spectrum. We find, $`K_0=4.33(1/\sqrt{g_s\alpha ^{}N})`$. However, it is easy to show that this eigenvalue does not correspond to a normalizable state. For large r, $`\sqrt{\mathrm{det}g}e^{\frac{5}{2}r}`$ and thus, $$\sqrt{\mathrm{det}g}_0^2e^{(\frac{5}{2}2\sqrt{(\frac{92}{5}+4.3^2)}r})$$ (3.36) so the integral $`_0^{\mathrm{}}\sqrt{\mathrm{det}g}_0^2`$ does not converge. The issue we are confronted with now is to find a good regularization for this model. Let us note an important point. Sean Hartnoll pointed that one might think of taking the norm in flat space $`𝑑r\|\|^2`$. Indeed, the equation (3.30) seems to indicate that, but we used a norm obtained from on the ten dimensional curved background $`d^{10}x\|\|^2`$, and is this norm the one that is forcing us to some regularization. His comment is based on the fact that one should only worry about our fluctuations to have finite energy and according to the paper , this condition implies the finitness of the norm in flat space. If we use his proposal, our wave functions $`(r)`$ are normalizable. In this paper we choose to attach to the more conventional norm defined in the curved space. The differences and physical implications of each choice of norm will be investigated elsewhere.<sup>4</sup><sup>4</sup>4We thank Sean Hartnoll and Stathis Tompaidis for valuable dicussions and input regarding this paragraph. Coming back to the regularization we need to introduce, let us recall the reader that, as explained in section two, in addition to the solution given by (2.1)-(2.8) there is another solution to the BPS equations, that we aluded to as the ”abelian” solution. The abelian (singular) and non-abelian (non-singular) solutions are the same at infinity but the abelian solution is singular in the IR and thus is not a good dual to SYM, that has nothing singular. We propose that a good regularization for this model is to use the background presented in (2.1)-(2.8) up to the point where it becomes indistinguishable from the abelian one. After this point the two solutions are the same and neither of them is dual to SYM, but to a higher dimensional theory on which we are not interested here. Figure (2) shows a plot of the Schrödinger potential for the abelian and non abelian solutions. Our proposal is that the IR solution that captures the physics of N=1 SYM is valid only up to the region $`\mathrm{\Lambda }_{Ab}`$ where the potential becomes the same as the one of the abelian solution. The scale $`\mathrm{\Lambda }_{Ab}`$, measured in units of $`g_s\alpha ^{}N`$, is set by the vacuum expectation value of the dilaton $`\mathrm{\Phi }_0`$. Numerically we are not doing anything new; the choice of the right endpoint of integration is always arbitrary, decided to best fit the physical problem at hand. Using a generalized shooting technique and integrating up to the point where $`V_{Ab}=V_{NonAb}\pm 0.00001`$, we find an improved value for the WKB estimate, $`K_0=4.291(1/\sqrt{g_s\alpha ^{}N})`$. This is a numerically stable eigenvalue meaning that small changes in the initital guess or pushing the endpoint of integration further to the right do not afect the value obtained. But this eigenstate is not normalizable, to obtain a normalizable state we have to impose the regularization procedure proposed above. This will be explained in detail in the next section. It is worth emphasizing that this model produces a discrete spectrum even without any kind of regularization. ### 3.1 Understanding The Regularization Procedure Above we have proposed that the correct way of computing in this model is to do a computation with the non-singular solution and, for large values of the radial coordinate, substract the result obtained with the singular background. This, we proposed is calculating in the dual N=1 SYM theory. Let us get a better intuition of this sort of “regularization procedure”. In physical terms, this procedure is easy to understand and is just instructing us to do our computations only in the region that is of interest to $`N=1`$ SYM. Indeed, since the non-abelian (non-singular) solution (that captures the IR effects of the dual field theory) asymptotes to the abelian (singular) solution (that is dual to a higher dimensional field theory), what we are basically doing when explicitely computing is substract the result obtained with the non-singular background minus those obtained at large values of the radial coordinate. Basically, this boils down to computing only in the region that is dual to $`N=1`$ SYM (with KK impurities as explained above). This is not very different from the type of regularizations done, for example in the computation of Wilson loops ,, where the infinite mass of the non-dynamical quark was substracted, or, what is the same, the mass of an infinite string not feeling the effects of the background (to which the real string asymptotes) is substracted. This sort of regularization was also used in . In that paper, even when a hard cut-off was imposed for numerical convenience, one should think about it as the fact that the computation was done in the region of interest, where the probe brane that adds flavor to the quenched version of N=1 SQCD is different from the probe brane in the singular background. More recently this sort of regularization was used in , even when the model used in that paper is different from ours, we believe their regularization can be understood in the lines we wrote above. Some readers might object the following: if one computes in this way, everything will give a finite result, so in this case, all functions will be normalizable. This questioning is valid, so let us try to answer it; for this it will be convenient to resort on an example where an exact solution is known. Hence, it is instructive to analyze the solution studied in the paper . Indeed, in that paper, the authors realized that a fluctuation given by $$\delta g_{\mu \nu }=\delta \varphi =0,\delta F_3.F_3=0,\delta F_3=_4dA_2$$ (3.37) with no restrictions to the functional form of the two form, solves the eqs of motion, that are written in the Appendix B (B.1)-(B.8). This could be a massless glueball, but when computing the norm $$d^{10}x\sqrt{g}\|\delta F_3\|^2,$$ seems to diverge, thus ruling it out as an state in the strong coupling theory. If we apply our criteria to this case, one might worry that the norm computed above will give a finite result, thus leading to a massless glueball that one does not expect in this theory (contrary to the KS case ). If these quantities give a finite result, this will imply that an effect not expected (a massless glueball) shows up. So, to understand this, let us do the computation for the norm, and apply our regularization procedure $$\delta F_3𝑑re^{2h+2\varphi }$$ (3.38) the regularization proposed above, indicates that we do a computation like this $$\delta F_3=_0^{r_0}𝑑re^{2\varphi +2h}\|_{nonsingular}+_{r_0}^{\mathrm{}}𝑑re^{2\varphi +2h}\|_{nonsingular}_{r_0}^{\mathrm{}}𝑑re^{2\varphi +2h}\|_{singular}$$ (3.39) Where $`r_0`$ is a value of the radial coordinate where the functions $`e^{2h+2\varphi }`$ computed in the singular and non-singular backgrounds are very similar to some degree of precision that is arbitrarily fixed. Since both integrands have the same asymptotics, they should equally diverge at large values of the radial coordinate. Indeed, both integrals diverge at leading order in the same way, but contrary to what one might expect, the integrals differ in a divergent quantity (and many convergent terms), thus, the computation in (3.39) is divergent and the configuration in (3.37) is not a good state of our theory. It is important to observe that many of the test that the solution has passed (see the review articles ), still work with this regularization. We would like to stress that even what was done in was not technically correct (as we mentioned, they seem to have used the wrong fluctuated eqs), the hard cut-off that they introduced is doing the same job that the regularization that we proposed here. Nevertheless, we have to make clear some important differences with (apart from the fact that we use different eqs.). In the authors did not find a discrete spectrum before imposing their hard cut-off, while we found one before our regularization, notice that our potential $`VS(r)`$ does ‘confine’ wavefunctions. We need to appeal to the regularization, only to satisfy the normalizability condition of our discrete states (and this is because we are being conservative and adopting a curved space measure for our normalizations). So, the hard cut-off regulation is quite different from what we have done here. If one introduces a cut-off, together with some boundary conditions, all solutions to the Schroedinger eqs will be normalizable. In our case, things are more subtle, as we explained above. ## 4 Glueballs from type IIA perspective This sections uses the same methods developed in Section 3 (See Appendix B.2 for all the details) for the type IIA background explained in Section 2. We will not carry out a full numerical analysis like in Section 3, but we will leave the system set for this more complicated (fully numerical) problem. Let us study glueballs for the case of wrapped branes in type IIA string theory. As explained above, in this case, the relevant background consists in $`N`$ D6 branes wrapping a three cycle inside a CY3 fold. So, to summarize the things clearly, let us write the dilaton and the metric of the IIA solution in Einstein frame, $`ds_E^2=2e^{\varphi /6}(dx_{1,3}^2+dr^2+(b(r)^2+g(r)^2)(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)+a(r)^2(d\stackrel{~}{\theta }^2+\mathrm{sin}^2\stackrel{~}{\theta }d\stackrel{~}{\phi }^2)+`$ $`+2g(r)a(r)^2[\mathrm{cos}\widehat{\psi }(d\theta d\stackrel{~}{\theta }+\mathrm{sin}\theta \mathrm{sin}\stackrel{~}{\theta }d\phi d\stackrel{~}{\phi })+\mathrm{sin}\widehat{\psi }(\mathrm{sin}\theta d\stackrel{~}{\theta }d\phi sin\stackrel{~}{\theta }d\theta d\stackrel{~}{\phi })]+`$ $`+{\displaystyle \frac{f^2c^2}{f^2+c^2(1+g_3)^2}}(d\widehat{\psi }+\mathrm{cos}\theta d\phi \mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\phi })^2),`$ $`4e^{4/3\varphi }=f(r)^2+c(r)^2(1+g_3(r))^2.`$ (4.40) The field strength $`F_2`$ reads, $$F_2=k^{}(r)dr(d\widehat{\psi }+\mathrm{cos}\theta d\phi \mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\phi })(k(r)+1)\mathrm{sin}\theta d\theta d\phi +(k(r)1)\mathrm{sin}\stackrel{~}{\theta }d\stackrel{~}{\theta }d\stackrel{~}{\phi },$$ (4.41) with $$k(r)=\frac{f^2c^2(1g_3^2)}{f^2+c^2(1+g_3)^2}.$$ In the following, we will work with this IIA set-up and because of the many similarities with the IIB model studied in the previous section, we will use the same approach. It is interesting to mention that if we can find glueballs in IIA, they should have an expression in M theory purely in terms of a metric fluctuation. Let us start by finding the dynamics of the fluctuations. The Lagrangian of IIA is <sup>5</sup><sup>5</sup>5notice that we take $`g_s=1`$ in this section $`L=\sqrt{g}\left(R{\displaystyle \frac{1}{2}}(\varphi )^2{\displaystyle \frac{1}{4}}e^{3/2\varphi }F_2^2{\displaystyle \frac{1}{12}}e^\varphi H_3^2{\displaystyle \frac{1}{48}}e^{\varphi /2}\widehat{F}_4^2\right)+{\displaystyle \frac{1}{2}}B_2F_4F_4,`$ $`\widehat{F}_4=dC_3H_3A_1,F_2=dA_1,H_3=dB_2,`$ (4.42) now, let us focus on the configurations that are of our interest, that is those where only the fields $`\varphi ,g_{\mu \nu },A_\mu `$ are turned on. The eqs of motion in this case are $`R_{\mu \nu }={\displaystyle \frac{1}{2}}_\mu \varphi _\nu \varphi +{\displaystyle \frac{1}{2}}e^{3/2\varphi }(F_\mu F_\nu ^{}{\displaystyle \frac{1}{16}}g_{\mu \nu }F^2),`$ $`R={\displaystyle \frac{1}{2}}(\varphi )^2+{\displaystyle \frac{3}{16}}e^{3/2\varphi }F^2,^2\varphi ={\displaystyle \frac{3}{8}}e^{3/2\varphi }F^2,`$ $`_\mu \left(\sqrt{g}e^{3/2\varphi }F^{\mu \nu }\right)=0,_{[\mu }F_{\nu \rho ]}=0`$ (4.43) Now, let us study the variations of these eqs. Under a fluctuation in all the relevant fields $`\varphi =\varphi +ϵ\delta \varphi ,g_{\mu \nu }=g_{\mu \nu }+ϵh_{\mu \nu }F_{\mu \nu }=F_{\mu \nu }+ϵ\delta F_{\mu \nu }.`$ (4.44) Again, we will propose a particular form for the metric fluctuation $$h_{ij}=\frac{\lambda }{5}hg_{ij},(i,j=x,t,r)h_{\alpha \beta }=h_{(\alpha \beta )}+\frac{h}{5}g_{\alpha \beta }.$$ (4.45) The same comments that we made regarding this decomposition, before eq. (3.25) are also pertinent here. We have denoted the trace of the internal part of the metric as $`h=h_a^a`$. We also propose an armonic expansion for the fields of a form similar to (3.26). $$h(x,y)=\mathrm{\Sigma }_Ih^I(x)Y^I(y),h_{(\alpha \beta )}=\mathrm{\Sigma }_Ib_{(\alpha \beta )}^I(x)Y^I(y),\delta \varphi (x,y)=\mathrm{\Sigma }_If^I(x)Y^I(y)$$ (4.46) Keeping only the s-wave fluctuations as before, we find after many computations (that are carefully spelled out in Appendix B), that like in IIB case, the many equations nicely combine and is sufficient to solve just one equation that reads, $$^2h^I+6g^{rr}_r\varphi _rh^I+\frac{e^{3/2\varphi }}{8}\left(31F^2+57(g^{\alpha \beta }F_\alpha F_\beta ^{}\frac{71}{57}g^{rr}F_rF_r^{})\right)h^I=0$$ (4.47) As we have done in the type IIB section, the numerics in this case can be studied. We will not do this here and we just want to point to the fact that even when this has to be done in a completely numerical way (since the metric functions are only numerically known), there is a nice feature of this IIA solution. The fact that the dilaton does not diverge, makes us believe that the regularization will not be necessary. This is left for future work. The point of this section was just to call the reader’s attention to this set of IIA models and show the analogy with the IIB treatment. ## 5 Summary, Conclusions and Future Directions Let us start by summarizing what we have done in this paper. First, we presented two Supergravity backgrounds (one in IIB and another in IIA) that are argued to be dual to $`N=1`$ SYM at low energies and we analyzed the spectrum in detail. Then, we initiated the study of glueball-like excitations in the strong coupling field theory as fluctuations of the Supergravity fields. We presented the equations to study the spectrum of glueballs and their excitations and analyzed them numerically. The type IIB case is analyzed in detail and we proposed a regularization procedure that might be useful in other computations involving these wrapped branes set-ups. Two appendixes present our computations in full detail. We find that unlike some IIB backgrounds previously studied, in the D5 wrapped on $`S^2`$ model not even the simplest scalar mode decouples from the rest of the fluctuations. Indeed, as we have shown, assuming only fluctuations of the dilaton leads to inconsistent equations. Therefore, the glueball $`0^{++}`$ in the IIB model we studied, is not dual to the dilaton, but to a mixture of dilaton and trace of the internal part of the metric. This goes in the same line as the papers , where the glueballs turned out to be mixings between different Supergravity modes. This mixing might persist for higher spin modes. The presence of a non-constant dilaton background seems to be the reason for the mixing of the fluctuations and this also appears in the model studied in . Another important point is that the potential found produces a discrete spectrum and a mass gap without any sort of cut-off, which seems to indicate that it is indeed capturing the physics of a confinig theory. As expected in a background with a linear dilaton, the states are not normalizable. Given that the UV completion of this theory is a little string theory it is not a surprise that a regularization (or substraction) procedure is needed. We present a proposal for this regularization which amounts to substracting the unwanted contribution of the UV regime. In the analysis of the type IIA background, we find that the scalar mode does not decouple from the rest of the fluctuations, indicating that, indeed, the non-constant dilaton in the background plays a role in producing this mixing. We do not perform the numerical analysis of the IIA background since the point of the paper is more to show a way of proceeding in these wrapped brane set-ups, that we believe is not exploited in the previous literature. Let us now discuss some future directions to follow and some work that should be interesting to do in detail, not discussed in this paper. First of all, as we mentioned in section 2, these models are contamined by the so called KK modes and we do not distinguish here if the glueballs we are obtaining are actually “made out” of KK modes composites (in which case they are not proper “glueballs” but hybrids composed of gluon, gluino and KK state) A good technical way to distinguish is to repeat the computation we are doing in Section 3, but for the case of the “dipole deformed” field theory (see for all details). Indeed, the idea in the paper is to make a $`SL(3,R)`$ deformation of the supergravity solution, that reflects on the dual field theory side on particular deformation that affects only the dynamics of the KK modes in the spectrum. Hence, the comparison of the eqs (3.27) and those in the appendix with those obtained in using the deformed metric can illuminate on what type of composition our ‘glueball’ has, if only glue and gluino, or if it is composed out of KK modes. Same could be repeated with deformations of the $`G_2`$ holonomy manifolds and the many examples already existing in the literature. It would be nice to check how our regularizing procedure works with the “flavor branes” introduced in when doing a dual to quenched SQCD. Indeed, there the idea was precisely the same, taking away from the computation the effects of the unwanted UV region. This can be understood by looking at the plots in figures 2 and 3 of . It should also be of interest, to study the glueballs in the type IIA model discussed in Section 4.1 of the paper . This solution is basically the same as the one we discussed in the IIB section, after some dualities. So, the interest of studying glueballs in this case is obviously to see if the same spectrum is obtained, analyze differences among eqs of motion, etc. Besides, it might have some interest to study the IIA case in more detail, not only numerically, as we pointed out above, but also from a $`G_2`$ holonomy perspective. Our dilaton-metric-gauge field fluctuations must combine in some way in a pure metric fluctuation in eleven dimensions. It should be nice to see how this works. Other models where it might have be interesting to apply our techniques (mainly the sort of manipulations explained in the Appendix B) is in models of non-supersymmetric duality. There is indeed one very clear model, that was studied in detail in the papers . One might also think about studying the ‘fermionic counterpart’ of what we have done in this paper, by fluctuating the fermionic fields around the bosonic background. On other respect, the comparison with the results from Lattice SYM should be done. There are some of these results in but we believe that the topic will evolve to allow better understanding and the contribution of this paper might be useful in the comparison with lattice results. See for example. Regarding this point, it should be interesting to study the spectrum of fermionic fluctuations (with fermionic fields vanishing in the background). This does not seem to have been exploited in the AdS/CFT literature, while other methods based on Veneziano-Yankielowicz and extensions, seem to give nice results . This situation might clearly improve with some study. The interest of studying glueballs goes beyond the simple fact of getting a discrete spectrum (that is by itself of enough interest). Indeed, glueballs play an important role in some recent advances regarding the study of Deep Inelastic and other types of Scattering using AdS/CFT techniques . The knowledge of glueballs masses and profiles in different models might help to extend the results in papers like to other ‘more realistic’ models. ## 6 Acknowledgments: We thank Richard Brower, José Edelstein, Nick Evans, Sean Hartnoll, Rafael Hernández, Oliver Jahn, Martin Kruczenski, Juan Martin Maldacena, Alfonso Ramallo, Angel Paredes, Chung-I-Tan, Pere Talavera and Stathis Tompaidis for discussions and comments that helped improving the presentation and interest of the results of this paper. Elena Cáceres would like to thank the Theory Group at the University of Texas at Austin for hospitality during the final stages of this work. This work was supported in part by the National Science Foundation under Grant No. PHY-0071512 and PHY-0455649, the US Navy, Office of Naval Research, Grant Nos. N00014-03-1-0639 and N00014-04-1-0336, Quantum Optics Initiative and by funds provided by the U.S.Department of Energy (DoE) under cooperative research agreement DF-FC02-94ER408818. Elena Cáceres is also supported by Mexico’s Council of Science and Technology, CONACyT, grant No.44840. Carlos Nuñez is a Pappalardo Fellow. ## Appendix A Appendix: Some Geometrical identities In the following, we list some geometrical identities that were used in the derivation of the fluctuated eqs (B.5)-(B.8) $`\delta R=_\mu _\nu h^{\mu \nu }^2h_\mu ^\mu R_{\mu \nu }h^{\mu \nu },`$ $`\delta R_{\mu \nu }={\displaystyle \frac{1}{2}}[_\alpha _\mu h_\nu ^\alpha +_\alpha _\nu h_\mu ^\alpha _\nu _\mu h_\alpha ^\alpha ^2h_{\mu \nu }],\sqrt{g+h}=\sqrt{g}(1+{\displaystyle \frac{ϵ}{2}}h_\mu ^\mu )+O(ϵ^2),`$ $`\delta \mathrm{\Gamma }_{\mu \nu }^\lambda ={\displaystyle \frac{1}{2}}g^{\alpha \lambda }(_\mu h_{\alpha \nu }+_\nu h_{\alpha \mu }_\alpha h_{\mu \nu })g^{\alpha \beta }\delta \mathrm{\Gamma }_{\alpha \beta }^r={\displaystyle \frac{1}{2}}g^{rr}(_rh_\mu ^\mu 2_\mu h_r^\mu )`$ $`\delta (^2\varphi )=^2\delta \varphi g^{ab}\delta \mathrm{\Gamma }_{ab}^\rho _\rho \varphi h^{ab}_a_b\varphi `$ (A.1) ## Appendix B Appendix: Derivation of the eqs in the IIB and IIA cases In this Appendix we fill in all the details that were left out in the computations that lead to eqs. (3.27) and (4.47) in sections three and four. ### B.1 Glueballs with the D5 branes solution Let us start with the type IIB solution. To study glueballs, we need the variation of the eqs of motion. In the type IIB case, for the solution of D5 branes wrapping $`S^2`$, the Einstein eqs for the metric, dilaton and three form read, $$R_{\mu \nu }=\frac{1}{2}_\mu \varphi _\nu \varphi +\frac{g_se^\varphi }{4}\left(F_\mu F_\nu ^{}\frac{1}{12}g_{\mu \nu }F_3^2\right)$$ (B.1) by contracting we get the Ricci scalar eq, $$R=\frac{1}{2}(_\mu \varphi )^2+\frac{g_se^\varphi }{24}F_3^2.$$ (B.2) The dilaton, Maxwell and Bianchi eqs. are, $$^2\varphi =\frac{g_se^\varphi }{12}F_3^2,_\mu \left(\sqrt{g}e^\varphi F^{\mu \nu \rho }\right)=0,_{[\mu }F_{\kappa \nu \rho ]}=0$$ (B.3) Now, let us study the fluctuations of these eqs. above. Let us assume that the background fields vary according to $$g_{\mu \nu }g_{\mu \nu }+ϵh_{\mu \nu },\varphi \varphi +ϵ\delta \varphi ,F_{\mu \nu \rho }F_{\mu \nu \rho }+ϵ\delta F_{\mu \nu \rho }$$ (B.4) So, keeping only linear order in the parameter $`ϵ`$, we get eqs for the fluctuated fields that read, for variation in the Ricci tensor eq. (3.21), $`{\displaystyle \frac{1}{2}}[_\alpha _\mu h_\nu ^\alpha +_\alpha _\nu h_\mu ^\alpha _\nu _\mu h_\alpha ^\alpha ^2h_{\mu \nu }]={\displaystyle \frac{1}{2}}[_\mu \varphi _\nu \delta \varphi +_\mu \delta \varphi _\nu \varphi ]`$ $`+{\displaystyle \frac{g_se^\varphi }{4}}[\delta \varphi F_\mu F_\nu ^{}+\delta F_\mu F_\nu ^{}+F_\mu \delta F_\nu ^{}2h^{ab}F_{\mu a}F_{\nu b}^{}]`$ $`{\displaystyle \frac{g_se^\varphi }{48}}[g_{\mu \nu }(2F_3\delta F_33h^{ab}F_aF_b^{}+\delta \varphi F_3^2)+h_{\mu \nu }F_3^2]`$ (B.5) and for the Ricci scalar eq.(B.2) $`_\mu _\nu h^{\mu \nu }^2h_\mu ^\mu R_{\mu \nu }h^{\mu \nu }=g^{\mu \nu }_\mu \varphi _\nu \delta \varphi {\displaystyle \frac{h^{\mu \nu }}{2}}_\mu \varphi _\nu \varphi +`$ $`{\displaystyle \frac{g_se^\varphi }{24}}[2F_3\delta F_33h^{kl}F_kF_l^{}+\delta \varphi F_3^2],`$ (B.6) For the fluctuated dilaton eq. we have, $$^2\delta \varphi h^{\mu \nu }_\mu _\nu \varphi \frac{1}{2}g^{rr}_r\varphi (2^\mu h_{r\mu }_rh_\mu ^\mu )\frac{g_se^\varphi }{12}(\delta \varphi F^2+2F\delta F3h^{\mu \nu }F_\mu F_\nu ^{})=0$$ (B.7) and for the fluctuation of the Maxwell eq. and Bianchi identity, $$_\mu [\sqrt{g}e^\varphi ((\frac{1}{2}h_\rho ^\rho +\delta \varphi )F^{\mu \nu \alpha }+\delta F^{\mu \nu \alpha }h^{\alpha c}F_c^{\mu \nu }+h^{\nu c}F_c^{\mu \alpha }h^{\mu c}F_c^{\nu \alpha })]=0,d\delta F=0$$ (B.8) Notice that the Ricci scalar eq. (B.6) can be obtained by contracting the Ricci tensor eq. (B.5) and substracting $`h^{\mu \nu }R_{\mu \nu }`$, so, in the following we will work with a contracted version of (B.5). In the derivation of these eqs, we have used the geometrical identities reviewed in Appendix A. Now, let us assume a fluctuation for the metric of the form $`h_{ij}={\displaystyle \frac{\lambda }{5}}g_{ij}h_\alpha ^\alpha ,(i,j=\stackrel{}{x},t,r)`$ $`h_{\alpha \beta }=h_{(\alpha \beta )}+{\displaystyle \frac{1}{5}}g_{\alpha \beta }h_a^a,(\alpha \beta )=(\theta ,\stackrel{~}{\theta },\phi ,\stackrel{~}{\phi },\psi )`$ (B.9) Where $`\lambda `$ is a constant and we have performed a shift in the $`h_{\alpha \beta }`$, this shift is proportional to the trace of the internal part of the metric $`h_\alpha ^\alpha `$. Note that the trace of the metric fluctuation over all the space is, $`h_\mu ^\mu =(\lambda +1)h_\alpha ^\alpha `$. Now, let us study the form of the fluctuation of the eq. that is obtained by contracting (B.5) with the background metric $`g^{\mu \nu }\delta R_{\mu \nu }`$, $$_\mu _\nu h^{\mu \nu }^2h_\mu ^\mu =^r\delta \varphi _r\varphi +\frac{g_se^\varphi }{24}[(\delta \varphi \frac{1}{2}g^{\mu \nu }h_{\mu \nu })F_3^2+3h^{\rho \sigma }F_\rho F_\sigma ^{}+2\delta F_3F_3]$$ (B.10) and using (B.9), $`{\displaystyle \frac{1}{5}}\lambda _x^2h_\alpha ^\alpha (\lambda +1)(_x^2+_y^2)h_\alpha ^\alpha g^{rr}_r\delta \varphi _r\varphi {\displaystyle \frac{g_se^\varphi }{24}}[(\delta \varphi {\displaystyle \frac{(\lambda +1)}{2}}h_\alpha ^\alpha )F_3^2`$ $`+{\displaystyle \frac{3}{5}}(\lambda g^{rr}F_rF_r^{}+g^{ab}F_aF_b^{})h_\alpha ^\alpha +2F_3\delta F_3+3h^{(ab)}F_aF_b^{}]=0.`$ (B.11) For the dilaton eq. (B.7) we will have, $`^2\delta \varphi +({\displaystyle \frac{3\lambda +5}{10}})g^{rr}_r\varphi _rh_\alpha ^\alpha {\displaystyle \frac{g_se^\varphi }{12}}[(\delta \varphi +{\displaystyle \frac{\lambda }{5}}h_\alpha ^\alpha )F^2+2F_3\delta F_3`$ $`{\displaystyle \frac{3}{5}}h_\alpha ^\alpha (\lambda g^{rr}F_rF_r^{}+g^{ab}F_aF_b^{})4h^{(\alpha \beta )}F_\alpha F_\beta ^{}]`$ (B.12) We have used the eq. of motion for the dilaton in the background eq.(3.23) above. The next step is to introduce an expansion in harmonics for each of the fields $$h_\alpha ^\alpha (x,y)=\mathrm{\Sigma }_I\mathrm{\Pi }^I(x)Y^I(y),h_{(\alpha \beta )}=\mathrm{\Sigma }_Ib^I(x)Y_{(\alpha \beta )}^I(y),\delta \varphi (x,y)=\mathrm{\Sigma }_If^I(x)Y^I(y)$$ (B.13) In order to satisfy the Bianchi identity we write the fluctuation $`\delta F_3`$ in the form $$\delta F_{\mu \nu \rho }=_{[\mu }\delta A_{\nu \rho ]}$$ (B.14) We will show that it is possible to find a fluctuation $`\delta F_3`$ orthogonal to the background $`F_3`$, i.e $`\delta F_3F_3=0`$, that satisfies Maxwell’s equation. Let us write the different componentes of Maxwell’s equation for a general fluctuation $`\delta F_3=_{[\mu }\delta A_{\nu \rho ]}`$ without demanding yet orthogonality with the backgorund $`F_3`$. We get $`\nu ,\alpha =a,b`$ (angular) $`_r[\sqrt{g}e^\varphi (({\displaystyle \frac{1}{2}}h_\mu ^\mu +\delta \varphi )F^{rab}+^{[r}\delta A^{ab]}h^{bc}F_c^{ra}+h^{ac}F_c^{rb}h^{rc}F_c^{ab})]`$ $`+_\theta [\sqrt{g}e^\varphi (({\displaystyle \frac{1}{2}}h_\mu ^\mu +\delta \varphi )F^{\theta ab}+^{[\theta }\delta A^{ab]}h^{bc}F_c^{\theta a}+h^{ac}F_c^{\theta b}h^{\theta c}F_c^{ab})]+_x[\sqrt{g}e^\varphi (^{[x}A^{ab]})]=0`$ . (B.15) $`\nu ,\alpha =r,b`$ (r,angular) $`_\theta [\sqrt{g}e^\varphi (({\displaystyle \frac{1}{2}}h_\mu ^\mu +\delta \varphi )F^{\theta rb}+^{[\theta }\delta A^{rb]}h^{bc}F_c^{\theta r}+h^{rc}F_c^{\theta b}h^{\theta c}F_c^{rb})]+_x[\sqrt{g}e^\varphi (^{[x}\delta A^{rb]})]=0`$ . (B.16) $`\nu ,\alpha =yz`$ $`_r[\sqrt{g}e^\varphi (^{[r}\delta A^{yz]})]+_\theta [\sqrt{g}e^\varphi (^{[\theta }A^{yz]})]+_x[\sqrt{g}e^\varphi (^{[x}\delta A^{yz]})]=0.`$ (B.17) $`\nu ,\alpha =(r,z)`$ $`_\theta [\sqrt{g}e^\varphi (^{[\theta }A^{rz]})]+_x[\sqrt{g}e^\varphi (^{[x}\delta A^{rz]})]=0.`$ (B.18) Now demand $`\delta F_3F_3=0`$, thus $`\delta F_{rab}=\delta F_{\theta ab}=0`$. From equations (B.15) and (B.16) above it is clear that a fluctuation of the form $`_x[\sqrt{g}e^\varphi (^{[x}A^{ab]})]`$ $`=`$ $`_r[\sqrt{g}e^\varphi (({\displaystyle \frac{1}{2}}h_\mu ^\mu +\delta \varphi )F^{rab}h^{bc}F_c^{ra}+h^{ac}F_c^{rb}h^{rc}F_c^{ab})]`$ $`+`$ $`_\theta [\sqrt{g}e^\varphi (({\displaystyle \frac{1}{2}}h_\mu ^\mu +\delta \varphi )F^{\theta ab}h^{bc}F_c^{\theta a}+h^{ac}F_c^{\theta b}h^{\theta c}F_c^{ab})]`$ $`_x[\sqrt{g}e^\varphi (^{[x}\delta A^{rb]})`$ $`=`$ $`_\theta [\sqrt{g}e^\varphi (({\displaystyle \frac{1}{2}}h_\mu ^\mu +\delta \varphi )F^{\theta rb}h^{bc}F_c^{\theta r}+h^{rc}F_c^{\theta b}h^{\theta c}F_c^{rb})]`$ with the other components given by (B.17) and (B.18) will satisfy Maxwell’s equation, the Bianchi identity and is such that $`\delta F_3F_3=0`$. Therefore, we have only eqs. (B.11) and (B.12), that keeping only the S-wave in the expansion in harmonics (B.13) take the form, $`({\displaystyle \frac{4\lambda }{5}}+1)_x^2\mathrm{\Pi }(x)^rf(x)_r\mathrm{\Phi }{\displaystyle \frac{g_se^\varphi }{24}}[(f(x){\displaystyle \frac{(\lambda +1)}{2}}\mathrm{\Pi }(x))F_3^2`$ $`+{\displaystyle \frac{3}{5}}(\lambda g^{rr}F_rF_r^{}+g^{ab}F_aF_b^{})\mathrm{\Pi }(x)]=0`$ (B.20) and $`^2f(x)+({\displaystyle \frac{3\lambda }{10}}+{\displaystyle \frac{1}{2}})^r\mathrm{\Pi }(x)_r\mathrm{\Phi }{\displaystyle \frac{g_s}{12}}e^\varphi [(f(x)+{\displaystyle \frac{\lambda }{5}}\mathrm{\Pi }(x))F_3^2`$ $`{\displaystyle \frac{3}{5}}(g^{ab}F_aF_b^{}+\lambda g^{rr}F_rF_r^{})\mathrm{\Pi }(x)]=0`$ (B.21) Equations (B.20) and (B.21) look suggestively similar so we will first check if there is any value of $`\lambda `$ for which they are the same. Indeed, for the two equations to be equal we need to satisfy $`2({\displaystyle \frac{4\lambda }{5}}+1)\mathrm{\Pi }(x)=f(x),`$ $`2f(x)=({\displaystyle \frac{3\lambda }{10}}+{\displaystyle \frac{1}{2}})\mathrm{\Pi }(x),`$ $`(f(x){\displaystyle \frac{(\lambda +1)}{2}}\mathrm{\Pi }(x))=f(x)+{\displaystyle \frac{\lambda }{5}}\mathrm{\Pi }(x)`$ (B.22) It is easy to check that $`\lambda =\frac{35}{29}`$ does the job. The eq. we need to solve is, $$^2f(x)+2g^{rr}_rf(x)_r\varphi +\frac{g_s}{12}e^\varphi \left[\frac{5}{2}F_3^2+\frac{293}{10}(g^{ab}F_aF_b^{}\frac{35}{29}g^{rr}F_rF_r^{})\right]f(x)=0$$ (B.23) This is precisely the eq.(3.27) we wanted to obtain. ### B.2 Glueballs with the D6 branes solution The relevant background consists in $`N`$ D6 branes wrapping a three cycle inside a CY3 fold. So,the metric of our IIA solution in Einstein frame and the dilaton where given in (4.40), and the Maxwell Field strength was, $$F_2=k^{}(r)dr(d\widehat{\psi }+\mathrm{cos}\theta d\phi \mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\phi })(k(r)+1)\mathrm{sin}\theta d\theta d\phi +(k(r)1)\mathrm{sin}\stackrel{~}{\theta }d\stackrel{~}{\theta }d\stackrel{~}{\phi },$$ (B.24) with $$k(r)=\frac{f^2c^2(1g_3^2)}{f^2+c^2(1+g_3)^2}.$$ Because of the many similarities with the IIB model studied in previous sections, we will use the same approach. Let us first study fluctuations. The Lagrangian of IIA is $`L=\sqrt{g}\left(R{\displaystyle \frac{1}{2}}(\varphi )^2{\displaystyle \frac{1}{4}}e^{3/2\varphi }F_2^2{\displaystyle \frac{1}{12}}e^\varphi H_3^2{\displaystyle \frac{1}{48}}e^{\varphi /2}\widehat{F}_4^2\right)`$ $`+{\displaystyle \frac{1}{2}}B_2F_4F_4,`$ $`\widehat{F}_4=dC_3H_3A_1,F_2=dA_1,H_3=dB_2.`$ (B.25) now, let us focus on the configurations that are of our interest, that is those where only the fields $`\varphi ,g_{\mu \nu },A_\mu `$ are turned on. The eqs of motion in this case are <sup>6</sup><sup>6</sup>6Like in the main text of the paper, we take $`g_s=1`$ in this Appendix $`R_{\mu \nu }={\displaystyle \frac{1}{2}}_\mu \varphi _\nu \varphi +{\displaystyle \frac{1}{2}}e^{3/2\varphi }(F_\mu F_\nu ^{}{\displaystyle \frac{1}{16}}g_{\mu \nu }F^2),`$ $`R={\displaystyle \frac{1}{2}}(\varphi )^2+{\displaystyle \frac{3}{16}}e^{3/2\varphi }F^2,^2\varphi ={\displaystyle \frac{3}{8}}e^{3/2\varphi }F^2,`$ $`_\mu \left(\sqrt{g}e^{3/2\varphi }F^{\mu \nu }\right)=0,_{[\mu }F_{\nu \rho ]}=0`$ (B.26) Now, let us study the variations of these eqs. Under a fluctuation in all the relevant fields $`\varphi =\varphi +ϵ\delta \varphi ,g_{\mu \nu }=g_{\mu \nu }+ϵh_{\mu \nu }F_{\mu \nu }=F_{\mu \nu }+ϵ\delta F_{\mu \nu }`$ (B.27) we have to first order in the fluctuation parameter $`ϵ`$, $`\delta R_{\mu \nu }={\displaystyle \frac{1}{2}}(_\mu \varphi _\nu \delta \varphi +_\nu \varphi _\mu \delta \varphi )+{\displaystyle \frac{1}{2}}e^{3/2\varphi }({\displaystyle \frac{3}{2}}\delta \varphi (F_\mu F_\nu ^{}{\displaystyle \frac{1}{16}}g_{\mu \nu }F^2)+F_\mu \delta F_\nu ^{}+`$ $`F_\nu \delta F_\mu ^{}h^{ab}F_{\mu a}F_{\nu b}{\displaystyle \frac{1}{8}}g_{\mu \nu }(F\delta Fh^{\alpha \beta }F_\alpha F_\beta ^{}){\displaystyle \frac{1}{16}}h_{\mu \nu }F^2),`$ $`\delta R=g^{\mu \nu }_\mu \varphi _\nu \delta \varphi {\displaystyle \frac{h^{\mu \nu }}{2}}_\mu \varphi _\nu \varphi +{\displaystyle \frac{3}{16}}e^{3/2\varphi }\left({\displaystyle \frac{3}{2}}\delta \varphi F^2+2F\delta F2h^{\mu \nu }F_\mu F_\nu ^{}\right),`$ $`\delta (^2\varphi )={\displaystyle \frac{3}{8}}e^{3/2\varphi }({\displaystyle \frac{3}{2}}\delta \varphi F^2+2F\delta F2h^{\mu \nu }F_\mu F_\nu ^{}),`$ $`_\mu \left(\sqrt{g}e^{3/2\varphi }[(3\delta \varphi +h_\mu ^\mu ){\displaystyle \frac{F^{\mu \nu }}{2}}+\delta F^{\mu \nu }h_a^\nu F^{\mu a}h_a^\mu F^{\nu a}]\right)=0`$ (B.28) So, using the geometrical variations given in the Appendix A and putting all together, we end up with three eqs, one for the variation of the dilaton, one for the Ricci scalar and the Maxwell eq. One can see that the eq for the variation of the Ricci tensor is included in the Ricci scalar variation. So, we have $$_\mu _\nu h^{\mu \nu }^2h_\mu ^\mu =g^{\mu \nu }_\mu \varphi _\nu \delta \varphi +\frac{e^{3/2\varphi }}{8}\left(\frac{9}{4}\delta \varphi F^2+3F\delta F+h^{\mu \nu }F_\mu F_\nu ^{}\frac{h_\mu ^\mu }{4}F^2\right)$$ (B.29) $$^2\delta \varphi h^{\mu \nu }_\mu _\nu \varphi \frac{g^{rr}}{2}(2_\mu h_r^\mu _rh_\mu ^\mu )_r\varphi =\frac{3e^{3/2\varphi }}{8}\left(\frac{3}{2}\delta \varphi F^2+2F\delta F2h^{\mu \nu }F_\mu F_\nu ^{}\right)$$ (B.30) $$_\mu \left(\sqrt{g}e^{3/2\varphi }[(3\delta \varphi +h_k^k)\frac{F^{\mu \nu }}{2}+\delta F^{\mu \nu }h_a^\nu F^{\mu a}h_a^\mu F^{a\nu }]\right)=0$$ (B.31) Now, let us propose a fluctuation of the metric $`h_{\mu \nu }`$ of the form $$h_{ij}=\frac{\lambda }{5}hg_{ij},(i,j=x,t,r)h_{\alpha \beta }=h_{(\alpha \beta )}+\frac{h}{5}g_{\alpha \beta }$$ (B.32) Here, we denote the trace of the internal part of the metric as $`h=h_a^a`$. So, let us study the fluctuated eqs, we will have for the Ricci scalar, $`({\displaystyle \frac{(4\lambda +5)}{5}}_x^2+{\displaystyle \frac{(4+5\lambda )}{5}}_y^2)h+g^{rr}_r\varphi _r\delta \varphi +{\displaystyle \frac{e^{3/2\varphi }}{8}}([{\displaystyle \frac{9\delta \varphi (\lambda +1)h}{4}}]F^2+3F\delta F`$ $`+{\displaystyle \frac{h}{5}}(g^{\alpha \beta }F_\alpha F_\beta ^{}+\lambda g^{rr}F_rF_r^{}))g^{(\alpha \beta )}(_\alpha _\beta h+F_\alpha F_\beta ^{})=0`$ (B.33) and for the dilaton, after the eq. of motion for the background dilaton field has been used we will have, $`^2\delta \varphi +({\displaystyle \frac{3\lambda +5}{10}})g^{rr}_r\varphi _rh{\displaystyle \frac{e^{3/2\varphi }}{8}}([{\displaystyle \frac{9}{2}}\delta \varphi +{\displaystyle \frac{3\lambda }{5}}h]F^2+6F\delta F{\displaystyle \frac{6}{5}}h(g^{\alpha \beta }F_\alpha F_\beta ^{}+\lambda g^{rr}F_rF_r^{})`$ $`{\displaystyle \frac{3}{4}}g^{(\alpha \beta )}F_\alpha F_\beta ^{})=0`$ (B.34) and the Maxwell eq. $`_\mu (\sqrt{g}e^{3/2\varphi }[(3\delta \varphi +h){\displaystyle \frac{F^{\mu \nu }}{2}}+\delta F^{\mu \nu }{\displaystyle \frac{\lambda }{5}}g^{\nu r}F_r^\mu g^{\nu a}hF_a^\mu g^{(\nu a)}F_a^\mu +`$ $`{\displaystyle \frac{\lambda }{5}}g^{\mu r}hF_r^\nu g^{\mu a}hF_a^\nu +g^{(\mu a)}F_a^\nu ])=0`$ (B.35) Now, let us follow an analysis very similar to the one we have done for the Type IIB case. First, we will propose an expansion of the form (B.13). $$h(x,y)=\mathrm{\Sigma }_Ih^I(x)Y^I(y),h_{(\alpha \beta )}=\mathrm{\Sigma }_Ib^I(x)Y_{(\alpha \beta )}^I(y),\delta \varphi (x,y)=\mathrm{\Sigma }_If^I(x)Y^I(y)$$ (B.36) Then, let us impose that both eqs are indeed the same. For this to happen, we will also need that the fluctuation of the Maxwell field is orthogonal to the field itself, that is $$\delta F_{\mu \nu }=_{[\mu }\delta A_{\nu ]},F_2\delta F_2=0.$$ (the first part is to automatically solve the Bianchi identity) and we will also concentrate on the $`s`$-wave, that is all modes with higher harmonics in the spheres will not be considered, same for $`g^{(\alpha \beta )}`$ that contains higher harmonics. As before, the armonic decomposition is such that $`_y^2h`$ is composed of higher harmonics. Then dividing eq (B.34) by 6, we have that the following equalities have to be satisfied $$f^I=6(\frac{4\lambda +5}{5})h^I=(\frac{3\lambda +5}{60})h^I,\mathrm{\hspace{0.33em}\hspace{0.33em}3}f^I=(\frac{\lambda +1}{4}\frac{\lambda }{10})h^I,$$ (B.37) for both eqs will be equal. Indeed, we can see that they are solved by $$\lambda =\frac{71}{57},\mathrm{\hspace{0.33em}\hspace{0.33em}95}f=2h,$$ and that both eqs. (B.33) and (B.34) actually read, $$^2h^I+6g^{rr}_r\varphi _rh^I+\frac{e^{3/2\varphi }}{8}\left(31F^2+57(g^{\alpha \beta }F_\alpha F_\beta ^{}\frac{71}{57}g^{rr}F_rF_r^{})\right)h^I=0$$ (B.38) Regarding the Maxwell eq, we can make an argument similar to the one in the IIB case .
warning/0506/math0506112.html	ar5iv	text	# Ginzburg-Weinstein via Gelfand-Zeitlin ## 1. Introduction and statement of results A theorem of Ginzburg-Weinstein states that for any compact Lie group $`K`$ with its standard Poisson structure, the dual Poisson Lie group $`K^{}`$ is Poisson diffeomorphic to the dual of the Lie algebra $`𝔨^{}`$, with the Kirillov Poisson structure. The result of does not, however, give a constructive way for obtaining such a diffeomorphism. For the case of the unitary group $`K=\mathrm{U}(n)`$, Flaschka-Ratiu (see also their preprint ) suggested the existence of a distinguished Ginzburg-Weinstein diffeomorphism, intertwining Gelfand-Zeitlin systems on $`𝔲(n)^{}`$ and $`\mathrm{U}(n)^{}`$, respectively. In this paper, we will give a proof of the Flaschka-Ratiu conjecture. The main result has the following ‘linear algebra’ implications, which may be stated with no reference to Poisson geometry. Let $`\mathrm{Sym}(n)`$ denote the space of real symmetric $`n\times n`$ matrices. For $`kn`$ let $`A^{(k)}\mathrm{Sym}(k)`$ denote the $`k`$th principal submatrix (upper left $`k\times k`$ corner) of $`A\mathrm{Sym}(n)`$, and $`\lambda _i^{(k)}(A)`$ its ordered set of eigenvalues, $`\lambda _1^{(k)}(A)\mathrm{}\lambda _k^{(k)}(A)`$. The map (1) $$\lambda :\mathrm{Sym}(n)^{\frac{n(n+1)}{2}},$$ taking $`A`$ to the collection of numbers $`\lambda _i^{(k)}(A)`$ for $`1ikn`$, is a continuous map called the Gelfand-Zeitlin map. Its image is the Gelfand-Zeitlin cone $`(n)`$, cut out by the ‘interlacing’ inequalities, (2) $$\lambda _i^{(k+1)}\lambda _i^{(k)}\lambda _{i+1}^{(k+1)},1ikn1.$$ Now let $`\mathrm{Sym}^+(n)\mathrm{Sym}(n)`$ denote the subset of positive definite symmetric matrices, and define a logarithmic Gelfand-Zeitlin map (3) $$\mu :\mathrm{Sym}^+(n)^{\frac{n(n+1)}{2}},$$ taking $`A`$ to the collection of numbers $`\mu _i^{(k)}(A)=\mathrm{log}(\lambda _i^{(k)}(A))`$. Then $`\mu `$ is a continuous map from $`\mathrm{Sym}^+(n)`$ onto $`(n)`$. ###### Theorem 1.1. There is a unique continuous map $`\psi :\mathrm{Sym}(n)\mathrm{SO}(n)`$, with $`\psi (0)=I`$, such that the map (4) $$\gamma =\mathrm{exp}\mathrm{Ad}_\psi :\mathrm{Sym}(n)\mathrm{Sym}^+(n),\mathrm{Ad}_\psi (A)\mathrm{Ad}_{\psi (A)}A$$ intertwines the Gelfand-Zeitlin maps $`\lambda `$ and $`\mu `$. In fact, $`\psi `$ is smooth and $`\gamma `$ is a diffeomorphism. Remark. For a general real semi-simple Lie group $`G`$ with Cartan decomposition $`G=KP`$, Duistermaat proved the existence of a smooth map $`\psi :𝔭K`$ such that the map $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi :𝔭P`$ intertwines the ‘diagonal projection’ with the ‘Iwasawa projection’. Theorem 1.1 gives canonical maps with this property for the case $`G=\mathrm{SL}(n,)`$. Example. The case $`n=2`$ can be worked out by hand (see also \[13, Example 3.27\]). Even in this case, smoothness of the map $`\gamma `$ is not entirely obvious. Since $`\gamma (A+tI)=e^t\gamma (A)`$, it is enough to consider trace-free matrices, $$A=\left(\begin{array}{cc}a& b\\ b& a\end{array}\right).$$ The matrix $`A`$ has Gelfand-Zeitlin variables $$\lambda _1^{(2)}(A)=r,\lambda _2^{(2)}(A)=r,\lambda _1^{(1)}(A)=a$$ with $`r:=\sqrt{a^2+b^2}`$. Hence, the matrix $`\gamma (A)`$ should have eigenvalues $`e^r,e^r`$ and upper left entry $`e^a`$. This gives $$\gamma (A)=\left(\begin{array}{cc}\stackrel{~}{a}& \stackrel{~}{b}\\ \stackrel{~}{b}& \stackrel{~}{c}\end{array}\right)$$ with $$\stackrel{~}{a}=e^a,\stackrel{~}{b}=\pm \sqrt{2e^a\mathrm{cosh}(r)e^{2a}1},\stackrel{~}{c}=2\mathrm{cosh}(r)e^a.$$ To obtain a continuous map, one has to take the sign of $`\stackrel{~}{b}`$ equal to the sign of $`b`$. The matrix $`\psi (A)\mathrm{SO}(2)`$ is a rotation matrix by some angle $`\theta (A)`$. A calculation gives, $$\mathrm{cos}(2\theta (A))=\frac{a}{r}\pm \sqrt{1\left(\frac{e^a\mathrm{cosh}(r)}{\mathrm{sinh}(r)}\right)^2}.$$ One can consider similar questions for the space $`\mathrm{Herm}(n)`$ of complex Hermitian $`n\times n`$-matrices, and its subset $`\mathrm{Herm}^+(n)`$ of positive definite matrices. Define surjective maps $$\lambda :\mathrm{Herm}(n)(n),\mu :\mathrm{Herm}^+(n)(n)$$ in terms of eigenvalues of principal submatrices, as before. Let $$\mathrm{Herm}_0(n)=\lambda ^1(_0(n))$$ denote the subset where all of the eigenvalue inequalities (2) are strict. The $`k`$-torus $`T(k)U(k)`$ of diagonal matrices acts on $`\mathrm{Herm}_0(n)`$ as follows, (5) $$tA=\mathrm{Ad}_{U^1tU}A,tT(k),A\mathrm{Herm}_0(n).$$ Here $`U\mathrm{U}(k)U(n)`$ is a unitary matrix such that $`\mathrm{Ad}_UA^{(k)}`$ is diagonal, with entries $`\lambda _1^{(k)},\mathrm{},\lambda _k^{(k)}`$. The action is well-defined since $`U^1tU`$ does not depend on the choice of $`U`$, and preserves the Gelfand-Zeitlin map (1). The actions of the various $`T(k)`$’s commute, hence they define an action of the Gelfand-Zeitlin torus $$T(n1)\times \mathrm{}\times T(1)\mathrm{U}(1)^{(n1)n/2}.$$ Here the torus $`T(n)`$ is excluded, since the action (5) is trivial for $`k=n`$. Let $`\mathrm{Herm}_0^+(n),\mathrm{Sym}_0(n)`$ and $`\mathrm{Sym}_0^+(n)`$ denote the intersections of $`\mathrm{Herm}_0(n)`$ with $`\mathrm{Herm}^+(n),\mathrm{Sym}(n)`$ and $`\mathrm{Sym}^+(n)`$. Thus $`\mathrm{Herm}_0^+(n)=\mu ^1(_0(n))`$. ###### Theorem 1.2. There is a unique continuous map $$\gamma :\mathrm{Herm}(n)\mathrm{Herm}^+(n)$$ with the following three properties: 1. $`\gamma `$ intertwines the Gelfand-Zeitlin maps: $`\mu \gamma =\lambda `$. 2. $`\gamma `$ intertwines the Gelfand-Zeitlin torus actions on $`\mathrm{Herm}_0(n)`$ and $`\mathrm{Herm}_0^+(n)`$. 3. For any connected component $`S`$ of $`\mathrm{Sym}_0(n)\mathrm{Herm}(n)`$, $`\gamma (S)S`$. In fact, $`\gamma `$ is a diffeomorphism, and has the following additional properties: 1. $`\gamma `$ is equivariant for the conjugation action of $`T(n)\mathrm{U}(n)`$, 2. $`\gamma (A+uI)=e^u\gamma (A)`$. 3. $`\gamma (\overline{A})=\overline{\gamma (A)}`$ (where the bar denotes complex conjugation). 4. For $`kn`$, the following diagram commutes: $$\begin{array}{ccccc}\mathrm{Herm}(n)& & \mathrm{Herm}(k)& & \mathrm{Herm}(n)\\ \gamma & & \gamma & & \gamma & & \\ \mathrm{Herm}^+(n)& & \mathrm{Herm}^+(k)& & \mathrm{Herm}^+(n)\end{array}$$ Here the left horizontal maps take a matrix to its $`k`$th principal submatrix, while the right horizontal maps are the obvious inclusions as the upper left corner, extended by $`0`$’s respectively $`1`$’s along the diagonal. Similar to the statement for real symmetric matrices, Theorem 1.1, the map $`\gamma `$ can be written in the form $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi `$ for a suitable map $`\psi :\mathrm{Herm}(n)\mathrm{SU}(n)`$. To fix the choice of $`\psi `$, we have to impose an equivariance condition under the Gelfand-Zeitlin torus action. Given $`A\mathrm{Herm}_0(n)`$, let $`U_kU(k)`$ be matrices diagonalizing $`A^{(k)}`$, and $`V_kU(k)`$ matrices diagonalizing $`\gamma (A)^{(k)}=\gamma (A^{(k)})`$. Then the Gelfand-Zeitlin action of $`t=(t_{n1},\mathrm{},t_1)T(n1)\times \mathrm{}\times T(1)`$ is given by $$tA=\mathrm{Ad}_{\chi (t,A)}A,t\gamma (A)=\mathrm{Ad}_{\stackrel{~}{\chi }(t,A)}\gamma (A)$$ where $$\chi (t,A)=U_1^1t_1U_1\mathrm{}U_{n1}^1t_{n1}U_{n1},\stackrel{~}{\chi }(t,A)=V_1^1t_1V_1\mathrm{}V_{n1}^1t_{n1}V_{n1}.$$ Note that $`\chi (t,A),\stackrel{~}{\chi }(t,A)`$ are independent of the choice of $`U_i,V_i`$. ###### Theorem 1.3. The map $`\psi :\mathrm{Sym}(n)\mathrm{SO}(n),\psi (0)=I`$ from Theorem 1.1 extends uniquely to a continuous (in fact, smooth) map $`\psi :\mathrm{Herm}(n)\mathrm{SU}(n)`$ with the equivariance property (6) $$\psi (tA)=\stackrel{~}{\chi }(t,A)\psi (A)\chi (t,A)^1$$ for all $`A\mathrm{Herm}_0(n),tT(n1)\times \mathrm{}\times T(1)`$. The map $`\gamma `$ from Theorem 1.2 is expressed in terms of $`\psi `$ as $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi `$. Furthermore, 1. $`\psi `$ is equivariant for the conjugation action of $`T(n)\mathrm{U}(n)`$. 2. $`\overline{\psi (A)}=\psi (\overline{A})`$. 3. For all $`kn`$, the following diagram commutes, $$\begin{array}{ccc}\mathrm{Herm}(k)& & \mathrm{Herm}(n)\\ \psi & & \psi & & \\ \mathrm{SU}(k)& & \mathrm{SU}(n)\end{array}$$ Note that the equivariance property (6) of $`\psi `$ implies the equivariance of $`\gamma `$: $$\gamma (tA)=\mathrm{exp}(\mathrm{Ad}_{\psi (tA)\chi (t,A)}A)=\mathrm{exp}(\mathrm{Ad}_{\stackrel{~}{\chi }(t,A)\psi (A)}A)=t\gamma (A).$$ Let us now place these results into the context of Poisson geometry. Let $`𝔲(n)`$ be the Lie algebra of $`\mathrm{U}(n)`$, consisting of skew-Hermitian matrices, and identify $$\mathrm{Herm}(n)𝔲(n)^{},$$ using the pairing $`A,\xi =2\mathrm{Im}(\mathrm{tr}A\xi )`$. Then $`\mathrm{Herm}(n)`$ inherits a Poisson structure from the Kirillov-Poisson structure on $`𝔲(n)^{}`$. It was proved by Guillemin-Sternberg in that the action of each $`T(k)`$ on $`\mathrm{Herm}_0(n)`$ is Hamiltonian, with moment map the corresponding Gelfand-Zeitlin variables, $`(\lambda _1^{(k)},\mathrm{},\lambda _k^{(k)})`$. On the other hand, the unitary group $`\mathrm{U}(n)`$ carries a standard structure as a Poisson Lie group, with dual Poisson Lie group $`\mathrm{U}(n)^{}`$ the group of complex upper triangular matrices with strictly positive diagonal entries. $`\mathrm{U}(n)^{}`$ may be identified with $`\mathrm{Herm}^+(n)`$, by the map taking the upper triangular matrix $`XU(n)^{}`$ to the positive Hermitian matrix $`(X^{}X)^{1/2}\mathrm{Herm}^+(n)`$. Flaschka-Ratiu proved that the $`T(k)`$ action on $`\mathrm{Herm}_0^+(n)`$ is Hamiltonian for the Poisson structure induced from $`\mathrm{U}(n)^{}`$, with moment map the logarithmic Gelfand-Zeitlin variables $`(\mu _1^{(k)},\mathrm{},\mu _k^{(k)})`$. ###### Theorem 1.4. The map $`\gamma :𝔲(n)^{}\mathrm{U}(n)^{}`$ described in Theorem 1.2 is a Poisson diffeomorphism. That is, for the group $`K=\mathrm{U}(n)`$ we have found a fairly explicit description of a Ginzburg-Weinstein diffeomorphism, in Gelfand-Zeitlin coordinates. By contrast, no coordinate expressions are known for the Ginzburg-Weinstein maps constructed in . Remark. A recent paper of Kostant-Wallach studies in detail the holomorphic (i.e., complexified) Gelfand-Zeitlin system, for the full space $`𝔤𝔩(n,)`$. It may be interesting to consider a nonlinear version of the holomorphic system, and to generalize our results to that setting. Acknowledgment. We would like to thank Henrique Bursztyn for helpful discussions. ## 2. Uniqueness of the map $`\gamma `$ In this Section, we construct a map $`\gamma `$ over the open dense subset $`\mathrm{Herm}_0(n)=\lambda ^1(_0(n))`$, satisfying all the properties listed in Theorem 1.2. The existence of a smooth extension to $`\mathrm{Herm}(n)`$ will be proved in the subsequent sections. We denote by $`T_{}(k)=T(k)\mathrm{O}(k)(_2)^k`$ the ‘real part’ of the torus. The action of the Gelfand-Zeitlin torus on $`\mathrm{Herm}_0(n)`$ restricts to an action of $`T_{}(n1)\times \mathrm{}T_{}(1)(_2)^{(n1)n/2}`$ on $`\mathrm{Sym}^+(n)`$. The following facts concerning the Gelfand-Zeitlin map are standard; we include the proof since we are not aware of a convenient reference. ###### Proposition 2.1. The restriction of the Gelfand-Zeitlin map to $`\mathrm{Herm}_0(n)`$ defines a principal bundle (7) $$\lambda :\mathrm{Herm}_0(n)_0(n)$$ with structure group the Gelfand-Zeitlin torus $`T(n1)\times \mathrm{}\times T(1)`$. It further restricts to a principal bundle (8) $$\lambda :\mathrm{Sym}_0(n)_0(n)$$ with structure group $`T_{}(n1)\times \mathrm{}\times T_{}(1)`$. Similarly for the restriction of the logarithmic Gelfand-Zeitlin map $`\mu :\mathrm{Herm}^+(n)(n)`$ to $`\mathrm{Herm}_0^+(n)`$ and $`\mathrm{Sym}_0^+(n)`$. ###### Proof. Consider the commutative diagram, (9) $$\begin{array}{ccc}\mathrm{Herm}_0(n)& & _0(n)\\ & & & & \\ \mathrm{Herm}_0(n1)& & _0(n1)\end{array}$$ where the horizontal maps are the Gelfand-Zeitlin maps, the left vertical map is $`AA^{(n1)}`$, and the right vertical map $`_0(n)_0(n1)`$ is the obvious projection, forgetting the variables $`\lambda _i^{(n)}`$. The Gelfand-Zeitlin map $`\mathrm{Herm}_0(n)_0(n)`$ factorizes as (10) $$\mathrm{Herm}_0(n)\mathrm{Herm}_0(n1)\times _{_0(n1)}_0(n)_0(n),$$ where the middle term is the fiber product. By induction, we may assume that the map $`\mathrm{Herm}_0(n1)_0(n1)`$, and hence the last map in (10), is a principal bundle for the action of the Gelfand-Zeitlin torus $`T(n2)\times \mathrm{}\times T(1)`$. Hence, it suffices to show that the first map in (10) is a principal $`T(n1)`$ bundle for the Gelfand-Zeitlin action. Thus, let $`\lambda _i^{(k)},1ikn`$ be the components of a given point $`\lambda _0(n)`$, and let $`A^{(n1)}\mathrm{Herm}_0(n1)`$ be a matrix with Gelfand-Zeitlin parameters $`\lambda _i^{(k)}`$ for $`1ikn1`$. Let us try to find $`b_1,\mathrm{},b_{n1}`$ and $`c`$ such that the matrix (11) $$A=\left(\begin{array}{cc}A^{(n1)}& b\\ b^{}& c\end{array}\right)$$ has eigenvalues $`\lambda _i^{(n)}`$. (Here $`b`$ denotes the $`(n1)\times 1`$-matrix with entries $`b_i`$.) Choose $`UU(n1)`$ such that the matrix $`\mathrm{\Lambda }^{(n1)}=UA^{(n1)}U^1`$ is diagonal, with entries $`\lambda _i^{(n1)}`$ down the diagonal. Then $$UAU^1=\left(\begin{array}{cc}\mathrm{\Lambda }^{(n1)}& \stackrel{~}{b}\\ \stackrel{~}{b}^{}& c\end{array}\right)$$ where $`\stackrel{~}{b}=Ub`$. (As before, we think of $`\mathrm{U}(k)`$ for $`kn`$ as a subgroup of $`\mathrm{U}(n)`$, using the inclusion as the upper left corner.) The characteristic polynomial $`det(AuI)`$ is therefore given by, $$det(AuI)=(cu)\underset{j}{}(\lambda _j^{(n1)}u)\underset{i}{}\|\stackrel{~}{b}_i\|^2\underset{ji}{}(\lambda _j^{(n1)}u).$$ Setting this equal to $`det(AuI)=_r(\lambda _r^{(n)}u)`$, and evaluating at $`u=\lambda _i^{(n1)}`$ and at $`u=\lambda _n^{(n)}`$, one finds $$\|\stackrel{~}{b}_i\|^2=\frac{_r(\lambda _r^{(n)}\lambda _i^{(n1)})}{_{ji}(\lambda _j^{(n1)}\lambda _i^{(n1)})},c=\lambda _n^{(n)}\underset{i}{}\frac{_{rn}(\lambda _r^{(n)}\lambda _i^{(n1)})}{_{ji}(\lambda _j^{(n1)}\lambda _i^{(n1)})}.$$ The eigenvalue inequalities ensure that the right hand side of the expression for $`\|\stackrel{~}{b}_i\|^2`$ is $`>0`$. This shows that the first map in (10) is onto. Furthermore, since $`c`$ is uniquely determined while $`\stackrel{~}{b}_i`$ are determined up to a phase, this map defines a principal $`T(n1)`$ bundle. Since left matrix multiplication of $`\stackrel{~}{b}`$ by an element of $`T(n1)`$ is exactly the Gelfand-Zeitlin action, the proof for $`\mathrm{Herm}_0(n)`$ is complete. The proof for $`\mathrm{Sym}_0(n)`$ is similar, considering only matrices with entries in $``$. The parallel statements for the map $`\mu `$ are a direct consequence of the statements for $`\lambda `$. ∎ ###### Lemma 2.2. There exists a unique continuous map $`\gamma :\mathrm{Herm}_0(n)\mathrm{Herm}_0^+(n)`$, satisfying (a)–(c) from Theorem 1.2. Furthermore, this map also has the Properties (d)–(g) from Theorem 1.2. ###### Proof. The choice of a connected component $`S\mathrm{Sym}_0(n)`$ defines a cross-section, hence a trivialization, of the principal bundle $`\lambda :\mathrm{Herm}_0(n)_0(n)`$. The intersection $$S^+=S\mathrm{Sym}_0^+(n)$$ is a connected component of $`\mathrm{Sym}_0^+(n)`$, which likewise trivializes the bundle $`\mu :\mathrm{Herm}_0^+(n)_0(n)`$. Thus, we obtain a unique map $`\gamma `$ satisfying (a)–(b), with $`\gamma (S)=S_+`$ for the given $`S`$. By equivariance, the property $`\gamma (S)=S_+`$ holds true for all components $`S\mathrm{Sym}_0(n)`$, which gives (c). We claim that Property (d) follows from (b). Indeed, the Gelfand-Zeitlin action of any element in $`t_kZ(U(k))T(k)T(n)`$ coincides with the conjugation action, since the functions $`\chi ,\stackrel{~}{\chi }`$ in (6) are simply $`\chi (t_k,A)=\stackrel{~}{\chi }(t_k,A)=t_k`$. The collection of these subgroups $`Z(U(k))\mathrm{U}(1)`$ of $`T(n)`$, together with the center $`Z(\mathrm{U}(n))`$ (which acts trivially) generate $`T(n)`$, proving the claim. Properties (e) and (f) follow from the uniqueness, since the maps $$Ae^u\gamma (A+uI),A\overline{\gamma (\overline{A})}$$ satisfy (a)–(c). Finally (g) holds by the commutativity of the diagram $$\begin{array}{ccccc}\mathrm{Herm}(n)& & \mathrm{Herm}(k)& & \mathrm{Herm}(n)\\ \lambda & & \lambda & & \lambda & & \\ (n)& & (k)& & (n)\end{array}$$ and of the similar diagram for map $`\mu `$. ∎ ###### Lemma 2.3. There is a continuous function $`\psi :\mathrm{Sym}_0(n)\mathrm{SO}(n)`$, with the property that the map $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi :\mathrm{Sym}_0(n)\mathrm{Sym}_0^+(n)`$ intertwines the Gelfand-Zeitlin maps. The map $`\psi `$ is uniquely determined by the additional condition $`\psi (uA)I`$ for $`u0`$. ###### Proof. We have seen above that there is a unique continuous map $`\gamma :\mathrm{Sym}_0(n)\mathrm{Sym}_0^+(n)`$ which intertwines the Gelfand-Zeitlin maps and satisfies $`\gamma (S)=S^+`$ for any component $`S\mathrm{Sym}_0(n)`$. Since $`\gamma (A)`$ and $`\mathrm{exp}(A)`$ have the same eigenvalues, and since $`S_0`$ is contractible, one can always choose a continuous map $`\psi :\mathrm{Sym}_0(n)\mathrm{SO}(n)`$ with $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi `$. Conversely, suppose $`\psi :\mathrm{Sym}_0(n)\mathrm{SO}(n)`$ is a continuous map, such that $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi `$ intertwines the Gelfand-Zeitlin maps. Suppose $`\psi (uA)I`$ for $`u0`$. We will show that (i) the map $`\gamma `$ has the property $`\gamma (S)=S^+`$ for any connected component $`S`$, and (ii) the map $`\psi `$ with these properties is unique. Proof of (i): It suffices to show that the restriction of $`\gamma `$ to $`\mathrm{Sym}_0(n)`$ is homotopic to the identity map of $`\mathrm{Sym}_0(n)`$. Define $$[0,1]\times \mathrm{Sym}_0(n)\mathrm{Sym}_0(n),(u,A)A_u=\{\begin{array}{cc}A\hfill & \text{ for }u=0\hfill \\ \frac{1}{u}(\gamma (uA)I)+uI\hfill & \text{ for }0<u1\hfill \end{array}$$ This is a well-defined continuous map since $$\underset{u0}{lim}\left(\frac{1}{u}(\gamma (uA)I)+uI\right)=\underset{u0}{lim}\left(\frac{1}{u}\left(\mathrm{exp}(\mathrm{Ad}_{\psi (uA)}A)I\right)\right)=A$$ Furthermore $`A_u\mathrm{Sym}_0(n)`$, since $`\mathrm{Sym}_0(n)`$ is invariant under scalar multiplication by nonzero numbers, as well as under addition of a scalar multiple of the identity matrix. Clearly $`A_1=\gamma (A)`$. Proof of (ii): Suppose $`\psi ^{}:\mathrm{Sym}_0(n)\mathrm{SO}(n)`$ is another map with $`\gamma (A)=\mathrm{exp}(\mathrm{Ad}_{\psi ^{}(A)}A)`$ and $`lim_{u0}\psi ^{}(uA)=I`$. Then $`\psi ^{}(A)=\psi (A)\chi (A)`$ where $`\chi (A)`$ centralizes $`A`$ and $`lim_{u0}\chi (uA)=I`$. Since the centralizer subgroup $`\mathrm{O}(n)_A`$ of any $`A\mathrm{Sym}_0(n)`$ is discrete, and $`\mathrm{O}(n)_A=\mathrm{O}(n)_{uA}`$ for $`u>0`$, we have $`\chi (A)=\chi (uA)\underset{u0}{\overset{}{}}I`$. Thus $`\chi (A)=I`$, proving uniqueness of $`\psi :\mathrm{Sym}_0(n)\mathrm{SO}(n)`$. ∎ Note that we have not yet shown that it is actually possible to satisfy the normalization condition $`lim_{u0}\psi (uA)=I`$. This can be proved ‘by hand’, but will in any case be automatic for the map constructed below (cf. Section 5.3). ###### Lemma 2.4. The map $`\psi :\mathrm{Sym}_0(n)\mathrm{SO}(n),lim_{u0}\psi (uA)=I`$ described in Lemma 2.3 admits a unique extension $`\psi :\mathrm{Herm}_0(n)\mathrm{SU}(n)`$ with the equivariance property (6). The composition $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi :\mathrm{Herm}_0(n)\mathrm{Herm}_0^+(n)`$ coincides with the map described in Lemma 2.2. Furthermore, this map also has the properties (a) – (c) described in Theorem 1.3. ###### Proof. By construction, the map $`\gamma :\mathrm{Sym}_0(n)\mathrm{Sym}_0^+(n)`$ has the equivariance property $`\gamma (tA)=t\gamma (A)`$ for all $`tT_{}(n1)\times \mathrm{}\times T_{}(1)`$. This implies the equivariance property (6) for the map $`\psi :\mathrm{Sym}_0(n)\mathrm{SO}(n)`$, using the uniqueness part of Lemma 2.3. Hence, $`\psi `$ admits a unique $`T(n1)\times \mathrm{}\times T(1)`$-equivariant extension to a map $`\mathrm{Herm}_0(n)\mathrm{SU}(n)`$, and the property $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi `$ follows by equivariance. Let us now check the additional properties from Theorem 1.3. (a) As mentioned above, the Gelfand-Zeitlin action of $`Z(\mathrm{U}(k))T(k)T(n)`$ for $`k<n`$ coincides with the action by conjugation. Hence, (6) gives $`\psi (\mathrm{Ad}_{t_k}A)=\mathrm{Ad}_{t_k}\psi (A)`$ for $`t_kZ(U(k))`$. Since the collection of these subgroups, together with $`Z(U(n))`$, generate $`T(n)`$ it follows that $`\psi `$ is $`T(n)`$-equivariant. (b) $`\overline{\psi (A)}=\psi (\overline{A})`$ follows from the uniqueness properties, since both $`\psi `$ and $`A\overline{\psi (\overline{A})}`$ are $`T(n1)\times \mathrm{}\times T(1)`$-equivariant extensions of the given map over $`\mathrm{Sym}_0(n)`$. (c) Let $`\psi ^{(k)}:\mathrm{Herm}(k)\mathrm{SU}(k)`$ denote the analogue of the map $`\psi `$, for given $`k<n`$. Since $`\psi `$ is equivariant for the conjugation by $`T(n)`$, it is in particular equivariant for the subgroup $`T(nk)`$ embedded as the lower right corner. Since $`\mathrm{Herm}_0(k)`$ (as the upper left corner) is fixed under this action, it follows that the restriction of $`\psi `$ takes values in $`\mathrm{S}(\mathrm{U}(k)\times T(nk))`$. Similarly, the restriction to $`\mathrm{Sym}_0(k)`$ takes values in $`\mathrm{S}(\mathrm{O}(k)\times T_{}(nk))`$. Since $`T_{}(nk)`$ is discrete, the property $`lim_{u0}\psi (uA)=I`$ implies that $`\psi \|\mathrm{Sym}_0(k)`$ must take values in $`\mathrm{SO}(k)`$. From the uniqueness properties, it therefore follows that it coincides with $`\psi ^{(k)}\|\mathrm{Sym}_0(k)`$. The more general statement $`\psi \|\mathrm{Herm}_0(k)=\psi ^{(k)}\|\mathrm{Herm}_0(k)`$ now follows by equivariance under the Gelfand-Zeitlin action of $`T(k1)\times \mathrm{}\times T(1)`$. ∎ To complete the proof of Theorems 1.1, 1.2, and 1.3, it suffices to find a smooth map $`\psi :\mathrm{Herm}(n)\mathrm{SU}(n),\psi (0)=I`$ with the following Properties: 1. $`\psi (\overline{A})=\overline{\psi (A)}`$, 2. $`\gamma =\mathrm{exp}\mathrm{Ad}_\psi `$ is a diffeomorphism intertwining the Gelfand-Zeitlin maps and the Gelfand-Zeitlin torus actions. 3. $`\psi `$ has the equivariance property (6). The construction of a map $`\psi `$ with these properties, using Poisson-geometric techniques, will be finished in Section 5.3. ## 3. Poisson-geometric techniques In this Section we discuss various tools that are needed for our construction of Ginzburg-Weinstein diffeomorphisms. ### 3.1. Bisections Suppose (12) $$𝒜:K\mathrm{Diff}(M)$$ is an action of a Lie group $`K`$ on a manifold $`M`$. We will often write $`k.x:=𝒜(k)(x)`$ for $`kK,xM`$. Consider the action groupoid $$K\times MM$$ with face maps $`_0(k,x)=x`$ and $`_1(k,x)=k.x`$. A bisection \[8, Chapter 15\] of $`K\times MM`$ is a submanifold $`NK\times M`$ such that both maps $`_i`$ restrict to diffeomorphisms $`NM`$. Any bisection has the form $`N=\{(x,\psi (x))\|xM\}`$ where $`\psi C^{\mathrm{}}(M,K)`$ is a map such that $$𝒜(\psi )(x):=𝒜(\psi (x))(x)$$ defines a diffeomorphism $`𝒜(\psi )\mathrm{Diff}(M)`$. Henceforth, we will refer to the map $`\psi `$ itself as a bisection. <sup>1</sup><sup>1</sup>1It can be shown that a smooth map $`\psi :MK`$ is a bisection, if and only if for all $`xM`$, the map $`k\psi (k.x)k`$ is a diffeomorphism of $`K`$. In this case, $`\psi ^1(x)=:h`$ is obtained as the unique solution of $`\psi (h.x)h=1`$. Let $`\mathrm{\Gamma }(M,K)C^{\mathrm{}}(M,K)`$ denote the set of bisections. The map (13) $$\mathrm{\Gamma }(M,K)\mathrm{Diff}(M),\psi 𝒜(\psi )$$ is a group homomorphism for the following product on $`\mathrm{\Gamma }(M,K)`$, $$(\psi _1\psi _2)(x)=\psi _1(𝒜(\psi _2)(x))\psi _2(x).$$ The inverse of a bisection $`\psi `$ for this product is given by $$\psi ^1(x)=\psi (𝒜(\psi )^1(x))^1.$$ The group homomorphism (13) extends the action (12) of $`K`$, and has kernel the bisections satisfying $`\psi (x)K_x`$ for all $`xM`$. For later reference we note the following easy fact: ###### Lemma 3.1. Suppose $`\psi \mathrm{\Gamma }(M,K)^K`$ is an equivariant bisection (that is, $`k\psi =\psi k`$ for all $`kK`$). Then $`\psi \varphi =\varphi \psi `$ for all bisections $`\varphi `$ satisfying $`𝒜(\psi )^{}\varphi =\varphi `$. Furthermore, $`(\varphi \psi )(x)=\varphi (x)\psi (x)`$. ###### Proof. Since $`𝒜(\psi )^{}\varphi =\varphi `$, the product $`(\varphi \psi )(x)`$ coincides with the pointwise product $`\varphi (x)\psi (x)`$. On the other hand, since $`\psi `$ is $`K`$-equivariant, $$(\psi \varphi )(x)=\psi (𝒜(\varphi )(x))\varphi (x)=\mathrm{Ad}_{\varphi (x)}(\psi (x))\varphi (x)=\varphi (x)\psi (x).$$ The Lie algebra $`\mathrm{\Gamma }(M,𝔨)`$ corresponding to $`\mathrm{\Gamma }(M,K)`$ may be described as follows. Let (14) $$𝔨𝔛(M),\xi 𝒜(\xi )$$ denote the infinitesimal generators of the $`K`$-action, i.e. $`𝒜(\xi )`$ is the vector field with flow <sup>2</sup><sup>2</sup>2In this paper, we follow the convention that the flow $`F_t`$ of a (possibly time dependent) vector field $`X_t`$ is defined in terms of its action on functions by $`(X_tf)(F_t^1(x))=\frac{}{t}f(F_t^1(x))`$. The Lie derivative $`L_{X_t}`$ on differential forms is then characterized by $`F_t^{}L_{X_t}=\frac{}{t}F_t^{}`$. $`F_t=𝒜(\mathrm{exp}(t\xi ))`$. Then (14) is a homomorphism of Lie algebras. For $`\beta C^{\mathrm{}}(M,𝔨)`$ let $`𝒜(\beta )𝔛(M)`$ be the vector field $`𝒜(\beta )(x)=𝒜(\beta (x))(x)`$. The map $`\beta 𝒜(\beta )`$ is a Lie algebra homomorphism for the ’action algebroid’ Lie bracket (15) $$[\beta _1,\beta _2](x)=[\beta _1(x),\beta _2(x)]+(L_{𝒜(\beta _1)}\beta _2)(x)(L_{𝒜(\beta _2)}\beta _1)(x).$$ on $`C^{\mathrm{}}(M,𝔨)`$. Let $`\mathrm{\Gamma }(M,𝔨)`$ denote the space $`C^{\mathrm{}}(M,𝔨)`$ with this Lie bracket. To see more clearly how $`\mathrm{\Gamma }(M,𝔨)`$ is the infinitesimal counterpart of $`\mathrm{\Gamma }(M,K)`$, it is useful to realize $`\mathrm{\Gamma }(M,K)`$ as a group of diffeomorphisms of $`K\times M`$. Define two commuting actions on $`K\times M`$, by setting $$\stackrel{~}{𝒜}(k)(h,x)=(hk^1,k.x),\stackrel{~}{𝒜}^{}(k)(h,x)=(kh,x).$$ Then the map $$\mathrm{\Gamma }(M,K)\mathrm{Diff}(K\times M),\psi \stackrel{~}{𝒜}(_0^{}\psi )$$ identifies $`\mathrm{\Gamma }(M,K)`$ with the group of diffeomorphism of $`K\times M`$ which commute with the action $`\stackrel{~}{𝒜}^{}`$ and preserve the $`\stackrel{~}{𝒜}`$-orbits. Similarly, $$\mathrm{\Gamma }(M,𝔨)𝔛(K\times M),\beta \stackrel{~}{𝒜}(_0^{}\beta )$$ identifies $`\mathrm{\Gamma }(M,𝔨)`$ with the Lie algebra of vector fields on $`K\times M`$ which are invariant under the action $`\stackrel{~}{𝒜}^{}`$ and are tangent to the $`\stackrel{~}{𝒜}`$-orbits. Let us now assume for simplicity that $`K`$ is compact. For any $`\beta \mathrm{\Gamma }(M,𝔨)`$, the vector field $`\stackrel{~}{𝒜}(_0^{}\beta )`$ is complete, since it is tangent to orbits. Hence, its time one flow exists, and defines an element of $`\mathrm{\Gamma }(M,K)`$. We have thus extended the exponential map $`\mathrm{exp}:𝔨K`$ to a map $$\mathrm{exp}:\mathrm{\Gamma }(M,𝔨)\mathrm{\Gamma }(M,K),$$ where $`\psi =\mathrm{exp}(\beta )`$ is the unique element such that $`\stackrel{~}{𝒜}(_0^{}\psi )`$ is the time one flow of $`\stackrel{~}{𝒜}(_0^{}\beta )`$. More generally, one can ‘integrate’ families of maps $`\beta _t\mathrm{\Gamma }(M,𝔨)`$ to families of bisections $`\psi _t`$, by viewing $`\beta _t`$ as a time dependent vector field $`\stackrel{~}{𝒜}(_0^{}\beta _t)`$ on $`K\times M`$ and identifying $`\psi _t`$ with the corresponding flow $`\stackrel{~}{F}_t`$ on $`K\times M`$. Equivalently, let $`F_t`$ be the flow of the vector field $`𝒜(\beta _t)`$ on $`M`$. Then $`F_t=𝒜(\psi _t)`$, where the bisection $`\psi _t\mathrm{\Gamma }(M,K)`$ is the solution of the ordinary differential equation on $`K`$, (16) $$\beta _t(F_t(x))=\frac{\psi _t(x)}{t}\psi _t(x)^1,\psi _0(x)=1.$$ ### 3.2. Gauge transformations of Poisson structures Let $`M`$ be a Poisson manifold, with Poisson bivector field $`\pi `$. The group of Poisson diffeomorphism of $`(M,\pi )`$ will be denoted $`\mathrm{Diff}_\pi (M)`$, and the group of Poisson vector fields by $`𝔛_\pi (M)`$. Let $`\sigma \mathrm{\Omega }^2(M)`$ be a closed 2-form, with the property that the bundle map (17) $$1+\sigma ^{\mathrm{}}\pi ^{\mathrm{}}:T^{}MT^{}M$$ is invertible everywhere on $`M`$. (Here $`\sigma ^{\mathrm{}}:TMT^{}M`$ and $`\pi ^{\mathrm{}}:T^{}MTM`$ are the bundle maps defined by a 2-form $`\sigma `$ and bivector field $`\pi `$, respectively.) Then the formula (18) $$\pi _\sigma ^{\mathrm{}}:=\pi ^{\mathrm{}}(1+\sigma ^{\mathrm{}}\pi ^{\mathrm{}})^1$$ defines a new Poisson structure $`\pi _\sigma `$ on $`M`$, called the gauge transformation of $`\pi `$ by $`\sigma `$ . The symplectic leaves of $`\pi _\sigma `$ coincide with those of $`\pi `$, while the leafwise symplectic forms change by the pull-back of $`\sigma `$. Gauge transformations of Poisson structures arise in the context of Hamiltonian group actions. A Poisson action $`𝒜:K\mathrm{Diff}_\pi (M)`$ is called Hamiltonian, if there exists a moment map $`\mathrm{\Phi }:M𝔨^{}`$, equivariant relative to the coadjoint action on $`𝔨^{}`$, such that the generating vector fields for the action are (19) $$𝒜(\xi )=\pi ^{\mathrm{}}\text{d}\mathrm{\Phi },\xi .$$ The moment map condition (19) shows in particular that Hamiltonian actions preserve the symplectic leaves. From the equivariance condition, it follows that $`\mathrm{\Phi }`$ is a Poisson map. Conversely, any Poisson map $`\mathrm{\Phi }:M𝔨^{}`$ defines a Lie algebra action by Equation (19). If $`K`$ is connected, and if the infinitesimal $`𝔨`$-action integrates to a $`K`$-action, then the latter is Hamiltonian with $`\mathrm{\Phi }`$ as its moment map. ###### Proposition 3.2. Let $`(M,\pi )`$ be a Hamiltonian $`K`$-manifold with moment map $`\mathrm{\Phi }`$. For any bisection $`\psi \mathrm{\Gamma }(M,K)`$ let (20) $$\sigma _\psi =\text{d}\mathrm{\Phi },(\psi ^1)^{}\theta ^L,$$ where $`\theta ^L\mathrm{\Omega }^1(K)𝔨`$ denotes the left-invariant Maurer-Cartan form. Then $`\sigma _\psi `$ defines a gauge transformation of $`\pi `$, and $$𝒜(\psi )_{}\pi =\pi _{\sigma _\psi }.$$ ###### Proof. Since it suffices to prove this identity leafwise, we may assume that $`\pi `$ is the inverse of a symplectic form $`\omega `$. The moment map condition (19) translates into $`\iota (𝒜(\xi ))\omega +\text{d}\mathrm{\Phi },\xi =0`$. We will show (21) $$𝒜(\psi ^1)^{}\omega =\omega +\sigma _\psi ,$$ thus in particular $`\omega +\sigma _\psi `$ is symplectic. One easily checks that the pull-back of $`\omega `$ under the map $`_1:K\times MM,(k,x)k.x`$ is $$_1^{}\omega =\omega \text{d}\mathrm{\Phi },\theta ^L\mathrm{\Omega }^2(K\times M).$$ Equation (21) follows, since $`𝒜(\psi ^1)`$ is a composition of $`_1`$ with the inclusion $`MK\times M,x(\psi ^1(x),x)`$. ∎ We collect some other formulas for the 2-form $`\sigma _\psi `$ which will become useful later. ###### Proposition 3.3. Let $`(M,\pi ,\mathrm{\Phi })`$ be as in Proposition 3.2. 1. For any bisection $`\psi \mathrm{\Gamma }(M,K)`$, $$𝒜(\psi )^{}\sigma _\psi =\text{d}\mathrm{\Phi },\psi ^{}\theta ^L.$$ 2. If $`\psi _1,\psi _2\mathrm{\Gamma }(M,K)`$ are bisections, $$\sigma _{\psi _1\psi _2}=\sigma _{\psi _1}+𝒜(\psi _1^1)^{}\sigma _{\psi _2}.$$ ###### Proof. (a) Using the equivariance of $`\mathrm{\Phi }`$ we have $$𝒜(\psi ^1)^{}\mathrm{\Phi },\psi ^{}\theta ^L=𝒜(\psi ^1)^{}\mathrm{\Phi },𝒜(\psi ^1)^{}\psi ^{}\theta ^L)=\mathrm{\Phi },\mathrm{Ad}_{\psi ^1}^1\left(𝒜(\psi ^1)^{}\psi ^{}\theta ^L\right).$$ But $`𝒜(\psi ^1)^{}\psi ^{}\theta ^L=(\psi ^1)^{}\theta ^R=\mathrm{Ad}_{\psi ^1}\left((\psi ^1)^{}\theta ^L\right)`$. (b) From the definition $`\psi _1\psi _2=(𝒜(\psi _2)^{}\psi _1)\psi _2`$ we obtain, $$(\psi _1\psi _2)^{}\theta ^L=\psi _2^{}\theta ^L+\mathrm{Ad}_{\psi _2()^1}(𝒜(\psi _2)^{}\psi _1^{}\theta ^L)$$ where $`\psi _2()^1`$ denotes the function $`x\psi _2(x)^1`$. (Not to be confused with $`\psi _2^1(x)`$.) Therefore, $$\text{d}\mathrm{\Phi },(\psi _1\psi _2)^{}\theta ^L=\text{d}\mathrm{\Phi },\psi _2^{}\theta ^L+𝒜(\psi _2)^{}\text{d}\mathrm{\Phi },\psi _1^{}\theta ^L.$$ Now apply $`𝒜((\psi _1\psi _2)^1)^{}=𝒜(\psi _1^1)^{}𝒜(\psi _2^1)^{}`$ to this result, and use (a). ∎ Proposition 3.3(b) shows in particular that $$\mathrm{\Gamma }_0(M,K)=\{\psi \mathrm{\Gamma }(M,K)\|\sigma _\psi =0\}$$ is a subgroup of the group of bisections. By Proposition 3.2, the homomorphism $`\mathrm{\Gamma }(M,K)\mathrm{Diff}(M),\psi 𝒜(\psi )`$ restricts to a group homomorphism, $$\mathrm{\Gamma }_0(M,K)\mathrm{Diff}_\pi (M).$$ ### 3.3. Moser’s method for Poisson manifolds Let $`(M,\pi )`$ be a Poisson manifold, and $`\sigma _t`$ a smooth family of closed 2-forms on $`M`$, with $`\sigma _0=0`$, such that $`1+\sigma _t^{\mathrm{}}\pi ^{\mathrm{}}`$ is invertible for all $`t`$. Consider the family of gauge-transformed Poisson structures, $`\pi _t=\pi _{\sigma _t}`$. Suppose (22) $$\frac{}{t}\sigma _t=\text{d}a_t$$ for a smooth family of 1-forms $`a_t\mathrm{\Omega }^1(M)`$, and define a time dependent Moser vector field $`v_t𝔛(M)`$ by $`v_t=\pi _t^{\mathrm{}}(a_t)`$. Assume that the time dependent vector field $`v_t`$ is complete (this is automatic if the symplectic leaves of $`M`$ are compact), and let $`F_t`$ be the flow with initial condition $`F_0=\mathrm{id}`$. One has , $$\pi _t=(F_t)_{}\pi .$$ The following alternative expression for the Moser vector field is useful: ###### Lemma 3.4. The Moser vector field is given by $`v_t=\pi ^{\mathrm{}}(b_t)`$ where $`b_t=a_t+\iota (v_t)\sigma _t`$. The 1-form $`b_t`$ satisfies, (23) $$\frac{}{t}\left(F_t^{}\sigma _t\right)=\text{d}(F_t^{}b_t)$$ where $`F_t`$ is the flow of $`v_t`$. ###### Proof. By definition $`v_t=\pi ^{\mathrm{}}(\stackrel{~}{b}_t)`$ where $`\stackrel{~}{b}_t=(1+\sigma _t^{\mathrm{}}\pi ^{\mathrm{}})^1a_t`$. The calculation $$a_t=(1+\sigma _t^{\mathrm{}}\pi ^{\mathrm{}})\stackrel{~}{b}_t=\stackrel{~}{b}_t\sigma _t^{\mathrm{}}v_t=\stackrel{~}{b}_t\iota (v_t)\sigma _t$$ shows $`\stackrel{~}{b}_t=b_t`$. From the definition of $`b_t`$ and the formula for $`\text{d}a_t`$, we find $$\text{d}b_t=\text{d}a_t+L(v_t)\sigma _t=\left(\frac{}{t}L(v_t)\right)\sigma _t=(F_t^1)^{}\frac{}{t}(F_t^{}\sigma _t).$$ We will refer to $`b_t`$ as the Moser 1-form. Note that for any given Poisson manifold $`(M,\pi )`$ the list of data $`v_t,F_t,a_t,b_t,\sigma _t,\pi _t`$ is determined by $`b_t`$ (and also by $`a_t`$). The following Proposition describes a situation where the twist flows $`𝒜(\psi _t)`$ from Section 3.1 can be viewed as Moser flows. ###### Proposition 3.5. Suppose $`K`$ is a compact Lie group, and $`(M,\pi )`$ is a Hamiltonian Poisson $`K`$-manifold with moment map $`\mathrm{\Phi }:M𝔨^{}`$. Let $`\beta _t\mathrm{\Gamma }(M,𝔨)`$ define (cf. (16)) the family of bisections $`\psi _t\mathrm{\Gamma }(M,K)`$ with $`\psi _0=1`$. Then the 2-form $`\sigma _t`$ determined by the Moser 1-form $$b_t=\text{d}\mathrm{\Phi },\beta _t$$ coincides with the form $`\sigma _{\psi _t}`$. Hence, the Moser flow $`F_t`$ coincides with $`𝒜(\psi _t)`$, and the gauge transformed Poisson structure $`\pi _t=\pi _{\sigma _t}`$ equals $`𝒜(\psi _t)_{}\pi `$. ###### Proof. By the moment map property, $`v_t=\pi ^{\mathrm{}}(b_t)=𝒜(\beta _t)`$, with flow $`F_t=𝒜(\psi _t)`$. We have to verify Equation (23). Observe (cf. (16)) $$\text{d}F_t^{}\beta _t=\text{d}\left(\frac{\psi _t}{t}\psi _t()^1\right)=\mathrm{Ad}_{\psi _t}\frac{}{t}\left(\psi _t^{}\theta ^L\right).$$ Since $`F_t^{}\mathrm{\Phi }=\mathrm{Ad}_{\psi _t()^1}^{}\mathrm{\Phi }`$ by equivariance of the moment map, this gives $$F_t^{}\mathrm{\Phi },\text{d}\beta _t=\mathrm{Ad}_{\psi _t()^1}^{}\mathrm{\Phi },\text{d}F_t^{}\beta =\mathrm{\Phi },\frac{}{t}(\psi _t^{}\theta ^L)=\frac{}{t}\mathrm{\Phi },\psi _t^{}\theta ^L.$$ Therefore, using Proposition 3.3(a), $$F_t^{}\text{d}b_t=\text{d}F_t^{}\mathrm{\Phi },\text{d}\beta _t=\frac{}{t}\text{d}\mathrm{\Phi },\psi _t^{}\theta ^L=\frac{}{t}F_t^{}\sigma _t.$$ ### 3.4. Stability of Poisson actions A well-known argument due to Palais shows that actions of compact Lie groups $`K`$ on compact manifolds $`M`$ are stable. That is, any deformation of such an action is obtained via conjugation by a family of diffeomorphisms of $`M`$. This result extends to the Poisson category: ###### Proposition 3.6 (Stability of Poisson actions of compact Lie groups). Let $`(M,\pi )`$ be a Poisson manifold, $`K`$ a compact Lie group, and $`𝒜_t:K\mathrm{Diff}_\pi (M)`$ a family of $`K`$-actions by Poisson diffeomorphism of $`M`$. Let $`w_t𝔛(M)`$ be the time dependent vector field, given in terms of its action on functions by (24) $$w_t=_K\text{d}k𝒜_t(k^1)^{}\frac{}{t}𝒜_t(k)^{}$$ where $`\text{d}k`$ denote the normalized Haar measure on $`K`$. Then $`w_t`$ is a Poisson vector field. If the flow $`F_t\mathrm{Diff}_\pi (M)`$ of $`w_t`$ exists (e.g. if $`M`$ is compact, or if the $`K`$-orbits are independent of $`t`$), then (25) $$𝒜_t(k)F_t=F_t𝒜_0(k),kK.$$ ###### Proof. The vector field $`w_t`$ given by (24) has the property, $$\frac{}{t}𝒜_t(k)^{}+w_t𝒜_t(k)^{}𝒜_t(k)^{}w_t=0.$$ Assuming that the flow $`F_t`$ of $`w_t`$ is defined, this integrates to, $$F_t^{}𝒜_t(k)^{}(F_t^1)^{}=𝒜_0(k)^{}.$$ which gives (25). Since $`𝒜_t(k)`$ are Poisson diffeomorphisms, each vector field $`w_t(k)=𝒜_t(k^1)^{}\frac{}{t}𝒜_t(k)^{}`$ is a Poisson vector field, and hence so is the $`K`$-average (24). ∎ Remark. Note that if the actions $`𝒜_t:K\mathrm{Diff}_\pi (M)`$ commute with another (fixed) action of a compact Lie group $`H`$, then the vector field $`w_t`$ and hence the diffeomorphisms $`F_t`$ are $`H`$-equivariant. ### 3.5. Poisson diffeomorphisms of $`𝔨^{}`$ Of particular importance is the case $`M=𝔨^{}`$, with $`𝒜:K\mathrm{Diff}_\pi (𝔨^{})`$ the coadjoint action. We begin with the following simple observation: ###### Lemma 3.7. For any compact Lie group $`K`$, the center of the group $`\mathrm{\Gamma }(𝔨^{},K)`$ of bisections is the subgroup $`\mathrm{\Gamma }(𝔨^{},K)^K`$ of equivariant bisections, and is contained in the kernel of the map $`\mathrm{\Gamma }(𝔨^{},K)\mathrm{Diff}(𝔨^{}),\psi 𝒜(\psi )`$. ###### Proof. Suppose $`\psi `$ is a $`K`$-equivariant bisection, i.e. $`k\psi =\psi k`$ for all $`kK`$. Equivalently, $`\psi (k.\mu )=\mathrm{Ad}_k\psi (\mu )`$ for all $`\mu 𝔨^{},kK`$. Specializing to $`kK_\mu `$, this shows that $`\psi (\mu )`$ is in the centralizer of $`K_\mu `$. Since $`K`$ is compact, this implies $`\psi (\mu )K_\mu `$. We have thus shown $`𝒜(\psi )=\mathrm{id}`$ for all $`\psi \mathrm{\Gamma }(𝔨^{},K)^K`$. Now suppose $`\psi ,\varphi \mathrm{\Gamma }(𝔨^{},K)`$, where $`\psi `$ is $`K`$-equivariant. Then $$(\psi \varphi )(\mu )=\psi (𝒜(\varphi )(\mu ))\varphi (\mu )=\mathrm{Ad}_{\varphi (\mu )}(\psi (\mu ))\varphi (\mu )=\varphi (\mu )\psi (\mu )=(\varphi \psi )(\mu )$$ (for the last equality, we used that $`𝒜(\psi )=\mathrm{id}`$). This shows that $`\mathrm{\Gamma }(𝔨^{},K)^K`$ is contained in the center of $`\mathrm{\Gamma }(𝔨^{},K)`$. The converse is obvious, since central elements commute in particular with elements of $`K`$. ∎ Remark. A similar statement holds for invariant open subsets of $`𝔨^{}`$, with the same proof. Consider $`k^{}`$ as a Hamiltonian $`K`$-space, with $`\mathrm{\Phi }`$ the identity map. The subgroup $`\mathrm{\Gamma }_0(𝔨^{},K)`$ of bisections $`\psi `$ with $`\sigma _\psi =0`$ is the group of Lagrangian bisections. (One can show that a bisection is Lagrangian if and only if its graph is a Lagrangian submanifold of the symplectic groupoid $`K\times 𝔨^{}=T^{}K`$.) The diffeomorphism $`𝒜(\psi )`$ defined by a Lagrangian bisection is a Poisson diffeomorphism preserving symplectic leaves. ###### Proposition 3.8. The kernel of the homomorphism (26) $$\mathrm{\Gamma }_0(𝔨^{},K)\mathrm{Diff}_\pi (𝔨^{}),\psi 𝒜(\psi )$$ is the group of invariant Lagrangian bisections $`\mathrm{\Gamma }_0(𝔨^{},K)^K`$, while its image is the normal subgroup of Poisson diffeomorphisms preserving symplectic leaves. ###### Proof. By Lemma 3.7 above, $`𝒜(\psi )=\mathrm{id}`$ for all $`\psi \mathrm{\Gamma }_0(𝔨^{},K)^K`$. Suppose conversely that $`\psi \mathrm{\Gamma }_0(𝔨^{},K)`$ is a Lagrangian bisection with $`𝒜(\psi )=\mathrm{id}`$. Then each $`\psi _t=r_t^{}\psi `$ generates the trivial action. In particular, this is true for the constant map $`\psi _0\psi (0)`$. Hence $`\psi (0)`$ is in the center of $`K`$. Replacing $`\psi `$ with $`\psi ^{}=\psi (0)^1\psi `$, we may assume $`\psi (0)=1`$. Let $`\beta _t=t^1r_t^{}\beta \mathrm{\Gamma }(𝔨^{},𝔨)^K`$ be the $`𝔨`$-valued functions generating $`\psi _t`$ (cf. (16)), and $`b_t=t^2r_t^{}b`$ the associated family of closed 1-forms. Since $`𝒜(\psi )=\mathrm{id}`$, the vector field $`v=\pi ^{\mathrm{}}(b)`$ is zero. Hence $`b`$ is $`K`$-invariant, and therefore $`\psi `$ is $`K`$-equivariant. Let $`F\mathrm{Diff}_\pi (𝔨^{})`$ be any Poisson diffeomorphism preserving symplectic leaves. (In particular, $`F(0)=0`$.) We have to show $`F=𝒜(\psi )`$ for some Lagrangian bisection $`\psi `$. Suppose first that $`F`$ is a linear Poisson diffeomorphism of $`𝔨^{}`$. Then $`F`$ is dual to a Lie algebra automorphism $`f\mathrm{Aut}(𝔨)`$. Since $`F`$ preserves orbits, the same is true for the map $`f`$. This implies that $`f`$ is an inner automorphism, $`f=\mathrm{Ad}_k`$ for some $`kK`$. Hence $`F=\mathrm{Ad}_k^{}=𝒜(k^1)`$ is given by a Lagrangian bisection $`\psi k^1`$. Consider now the general case. For $`t`$ let $`r_t:𝔨^{}𝔨^{}`$ denote the map $`r_t(\mu )=t\mu `$. Let $`F_t=r_{t^1}Fr_t`$ for $`t0`$, so that the limit for $`t0`$ is the linearization $`F_0=\text{d}_0F`$ of $`F`$ at the origin. Since $`(r_t)_{}\pi =t\pi `$, each $`F_t`$ is a Poisson diffeomorphism preserving leaves. By the linear case considered above, we may assume $`F_0=\mathrm{id}`$. Let $`v_t𝔛_\pi (M)`$ be the time-dependent vector field, given in terms of its action on functions by $`v_t=(F_t^1)^{}\frac{}{t}F_t^{}`$. Write $`v=v_1`$. Then $`v_t=t^1(r_{t^1})_{}v`$ for $`t0`$. The vector fields $`v_t`$ vanish to second order at $`0`$, since $`F_t(0)=0`$ and $`\text{d}_0F_t\mathrm{id}`$ for all $`t`$. In particular, $`v_0=0`$. We now use the well-known fact that a Poisson vector field on $`𝔨^{}`$ is Hamiltonian if and only if it is tangent to the symplectic leaves (which is automatic if $`𝔨`$ is semi-simple). This follows from the description of the first Poisson cohomology of $`𝔨^{}`$ (see e.g. ) $$H_\pi ^1(𝔨^{})(𝔨^{})^KC^{\mathrm{}}(𝔨^{})^K.$$ Hence, we may write $`v=\pi ^{\mathrm{}}(b)`$ for some exact 1-form $`b\mathrm{\Omega }^1(𝔨^{})`$. The 1-form $`b`$ can be normalized by requiring that its $`K`$-average be zero. (Note that exact 1-forms on $`𝔨^{}`$ generate the zero vector field if and only if they are $`K`$-invariant.) Letting $`b_t=t^2r_t^{}b`$, and denoting by $`\beta _t=t^1r_t^{}\beta \mathrm{\Gamma }(𝔨^{},𝔨)`$ the corresponding $`𝔨`$-valued functions, we get $$v_t=\pi ^{\mathrm{}}(b_t)=𝒜(\beta _t).$$ Let $`\psi _t\mathrm{\Gamma }(𝔨^{},K),\psi _0=1`$ be the family of bisections obtained by integrating $`\beta _t`$ (see (16)). We have $`\psi _t=r_t^{}\psi `$ with $`\psi =\psi _1`$. Since the 1-forms $`b_t`$ are closed, the corresponding 2-forms $`\sigma _t=\sigma _{\psi _t}`$ (cf. (23) and Proposition 3.5) vanish. That is, the bisections $`\psi _t`$ are Lagrangian. We have $`F_t=𝒜(\psi _t)`$ by construction, and in particular $`F=𝒜(\psi )`$. ∎ ###### Remark 3.9. Let $`(M,\pi )`$ be a Poisson manifold admitting a symplectic realization $`SM`$. In Bursztyn-Weinstein \[6, Section 5\], the Poisson diffeomorphisms of $`M`$ which are generated by Lagrangian bisections of $`S`$ are referred to as inner automorphisms of $`M`$. Proposition 3.8 characterizes the inner automorphisms for the case $`T^{}K𝔨^{}`$. ###### Proposition 3.10. Suppose $`\sigma \mathrm{\Omega }^2(𝔨^{})`$ is a closed 2-form, defining a gauge transformation of the Kirillov-Poisson structure $`\pi `$ on $`𝔨^{}`$. Then there exists a bisection $`\psi \mathrm{\Gamma }(𝔨^{},K)`$, $`\psi (0)=1`$, such that $`\sigma _\psi =\sigma `$. In particular, $`𝒜(\psi )_{}\pi =\pi _\sigma `$. $`\psi `$ is unique up to multiplication from the right by a Lagrangian bisection. If $`\sigma `$ is invariant under the action of $`HK`$, the bisection $`\psi `$ can be taken $`H`$-invariant. ###### Proof. The assumption on $`\sigma `$ means that the bundle map $`A=1+\sigma ^{\mathrm{}}\pi ^{\mathrm{}}`$ is invertible everywhere. Define a smooth family of closed 2-forms $`\sigma _t`$, by letting $`\sigma _0=0`$ and $`\sigma _t=t^1r_t^{}\sigma `$ for $`t0`$. Introduce the corresponding operators $$A_t=1+\sigma _t^{\mathrm{}}\pi ^{\mathrm{}}$$ on $`T^{}𝔨^{}`$, connecting $`A_1=A`$ with $`A_0=1`$. Using $`(r_t)_{}\pi =t\pi `$, one finds $`A_t=r_t^{}Ar_{t^1}^{}`$ for $`t0`$. Since $`A`$ is invertible, it follows that the operator $`A_t`$ is invertible for all $`t`$. Hence, each $`\sigma _t`$ defines a gauge transformation. Now let $`a_t`$ be the family of 1-forms, obtained by applying the homotopy operator to $`\frac{}{t}\sigma _t`$, and $`b_t`$ the corresponding family of Moser 1-forms. By Proposition 3.5, the bisections $`\psi _t`$ corresponding to $`b_t`$ satisfy $`\sigma _{\psi _t}=\sigma _t`$. Thus $`\psi =\psi _1`$ has the desired property $`\sigma _\psi =\sigma `$. Uniqueness of $`\psi `$ up to Lagrangian bisections follows directly from Proposition 3.3(b). If $`\sigma `$ is $`H`$-invariant, then the bisection $`\psi `$ just constructed is $`H`$-invariant. ∎ For any compact, connected Lie group $`K`$, we denote by $`Z(K)K`$ the identity component of the center, and by $`K^{ss}`$ its semi-simple part (commutator subgroup). Thus $`\widehat{K}=K^{ss}\times Z(K)K`$ is a finite covering of $`K`$, and $`𝔨=𝔨^{ss}𝔷(𝔨)`$ on the level of Lie algebras. ###### Proposition 3.11. Let $`K_1,K_2`$ be compact Lie groups, and suppose $`\mathrm{\Phi }:𝔨_2^{}𝔨_1^{}`$ is the moment map for a Hamiltonian action $`𝒜:K_1\mathrm{Diff}_\pi (𝔨_2^{})`$. Suppose that the composition of $`\mathrm{\Phi }`$ with the projection $`𝔨_1^{}𝔷(𝔨_1)^{}`$ is a linear map, $`𝔨_2^{}𝔷(𝔨_1)^{}`$. Then there exists a Lie algebra homomorphism $`\tau :𝔨_1𝔨_2`$ and a Lagrangian bisection $`\psi \mathrm{\Gamma }_0(𝔨_2^{},K_2)`$ such that $`\mathrm{\Phi }=\tau ^{}𝒜(\psi )`$. ###### Proof. Let us first of all observe that $`\mathrm{\Phi }(0)=0`$. Indeed, for the $`𝔷(𝔨_1)^{}`$-component of $`\mathrm{\Phi }`$ this follows by linearity, while for the $`(𝔨_1^{ss})^{}`$-component it follows since moment maps are equivariant by definition. For all $`t0`$, the scaled Poisson homomorphism $`\mathrm{\Phi }_t=r_t^1\mathrm{\Phi }r_t`$ is a moment map for the scaled action $`k𝒜_t(k)=r_t^1𝒜(k)r_t`$. Note that the $`𝔷(𝔨_1)^{}`$-component of $`\mathrm{\Phi }_t`$, and hence the $`Z(K_1)K_1`$-action, do not depend on $`t`$. The limit $`\mathrm{\Phi }_0`$ for $`t0`$ equals the linearization of $`\mathrm{\Phi }`$ at $`0`$, and is a moment map for the linearized action $`𝒜_0`$. By Proposition 3.6 and the subsequent Remark, there exists a $`Z(K_1)`$-equivariant Poisson diffeomorphism $`F\mathrm{Diff}_\pi (𝔨_2^{})`$ with $`𝒜(k)=F𝒜_0(k)F^1,kK_1`$. Since moment maps for semisimple Lie groups (in this case $`K_1^{ss}`$) are unique, and since the $`𝔷(𝔨_1)^{}`$-component of $`\mathrm{\Phi }`$ does not depend on $`t`$, this implies $`\mathrm{\Phi }=\mathrm{\Phi }_0F^1`$. By Proposition 3.8, $`F^1=𝒜(\psi )`$ for some Lagrangian bisection $`\psi `$. Since $`\mathrm{\Phi }_0`$ is a linear Poisson map, it is of the form $`\mathrm{\Phi }_0=\tau ^{}`$ for a Lie algebra homomorphism $`\tau ^{}`$. ∎ ## 4. Ginzburg-Weinstein diffeomorphisms The main result of this Section is Theorem 4.7, showing that Ginzburg-Weinstein diffeomorphisms can be arranged to be compatible with given Poisson Lie group homomorphisms. ### 4.1. Poisson Lie groups We briefly review Poisson Lie groups, referring to for more detailed information. A Poisson Lie group is a Lie group $`K`$, equipped with a Poisson structure $`\pi ^K`$ for which the product map is Poisson. The linearization of $`\pi ^K`$ at the group unit is a Lie algebra 1-cocycle $`\delta ^K:𝔨^2𝔨`$, with the property that the dual map $`(\delta ^K)^{}`$ defines a Lie bracket on $`𝔨^{}`$. Conversely, if $`K`$ is connected, the cobracket $`\delta ^K`$ determines $`\pi ^K`$. One refers to the Lie algebra $`𝔨`$ together with $`\delta ^K`$ as the tangent Lie bialgebra of the Poisson Lie group $`K`$. The dual Poisson Lie group $`K^{}`$ is the connected, simply connected Poisson Lie group with tangent Lie bialgebra $`𝔨^{}`$. If $`\pi ^K=0`$, the dual Poisson Lie group is simply $`𝔨^{}`$ with the Kirillov Poisson structure. A Poisson Lie group action of $`(K,\pi ^K)`$ on a Poisson manifold $`(M,\pi ^M)`$ is a $`K`$-action such that the action map $`_1:K\times MM,(k,x)k.x`$ is a Poisson map. A $`K^{}`$-valued moment map for such an action is a Poisson map $`\mathrm{\Psi }:MK^{}`$ such that the generating vector fields are given by (27) $$𝒜(\xi )=(\pi ^M)^{\mathrm{}}\mathrm{\Psi }^{}\theta ^R,\xi .$$ Here $`\theta ^R\mathrm{\Omega }^1(K^{})𝔨^{}`$ is the right-invariant Maurer-Cartan form on $`K^{}`$. Equation (27) reduces to the usual moment map condition (19) if $`K`$ carries the zero Poisson structure. According to Lu , any Poisson map $`\mathrm{\Psi }:MK^{}`$ defines an infinitesimal Poisson Lie group action via (27). In particular, the identity map of $`K^{}`$ defines an infinitesimal dressing action of $`K`$ on $`K^{}`$. In nice cases, it integrates to a global action of $`K`$. Compact Lie group $`K`$ carry a standard Poisson structure $`\pi ^K`$ structure, constructed by Lu and Weinstein in . Let $`G=K^{}`$ be the complexification of $`K`$, viewed as a real Lie group, and $`𝔤=𝔨^{}`$ its Lie algebra. Consider the Iwasawa decomposition $$𝔤=𝔨𝔞𝔫,G=KAN$$ relative to a choice of maximal torus $`TK`$ and fundamental Weyl chamber. That is, $`𝔞=\sqrt{1}𝔱`$ while $`𝔫`$ is the direct sum of root spaces for the positive roots. Let $`B(,)`$ be an invariant scalar product on $`𝔨`$, and $`B^{}(,)`$ its complexification. Then $`2\mathrm{Im}B^{}(,)`$ is an invariant scalar product on $`𝔤`$, and restricts to a non-degenerate pairing between $`𝔨`$ and the Lie algebra $`𝔞𝔫`$. In this way $`𝔨^{}𝔞𝔫`$ acquires a Lie algebra structure, making $`𝔨`$ into a Lie bialgebra. Thus $`K`$ is a Poisson Lie group, with $`K^{}=AN`$ the dual Poisson Lie group. The action of $`K`$ on $`G`$ by left multiplication descends to the dressing action $`𝒜_K^{}`$ on $`K^{}`$, viewed as a homogeneous space $`G/K`$. To analyze the dressing action, it is convenient to work with the Cartan decomposition (28) $$𝔤=𝔨𝔭,G=KP.$$ where $`𝔭=\sqrt{1}𝔨`$ and $`P=\mathrm{exp}𝔭`$. Recall that the exponential map $`\mathrm{exp}:𝔤G`$ restricts to a diffeomorphism $`𝔭P`$. Let $`e:𝔨^{}K^{}`$ be the diffeomorphism defined by the compositions, $$𝔨^{}𝔤/𝔨𝔭\stackrel{\mathrm{exp}}{}PG/KK^{}.$$ Then $`e`$ intertwines the coadjoint action $`𝒜_𝔨^{}`$ with the dressing action: $$e𝒜_𝔨^{}(k)=𝒜_K^{}(k)e.$$ Example. Let $`K=\mathrm{U}(n)`$, with maximal torus $`T=T(n)`$ and the usual choice of positive roots. Then $`G=\mathrm{GL}(n,)`$ (viewed as a real Lie group), $`N`$ are the upper triangular matrices with $`1`$’s down the diagonal, and $`A`$ are the diagonal matrices with positive entries. Hence $`K^{}=AN`$ are the upper triangular matrices with positive diagonal entries. Furthermore, $`𝔭=\mathrm{Herm}(n)`$ and $`P=\mathrm{Herm}^+(n)`$. The isomorphism $`K^{}P`$ is explicitly given by $`X(X^{}X)^{1/2}`$, and identifies the dressing action with the conjugation action on $`\mathrm{Herm}^+(n)`$. ### 4.2. Ginzburg-Weinstein diffeomorphisms Let $`K`$ be a compact Lie group with the standard Poisson structure, and consider the map $`e:𝔨^{}K^{}`$ constructed above. In , it was observed that the Poisson structure $`\pi _1^𝔨^{}=(e^1)_{}\pi ^K^{}`$ is gauge equivalent to the Kirillov-Poisson structure $`\pi _0^𝔨^{}=\pi ^𝔨^{}`$. ###### Theorem 4.1. There is a canonical $`T`$-invariant closed 2-form $`\sigma \mathrm{\Omega }^2(𝔨^{})`$, with the property $$(e^1)_{}\pi ^K^{}=\pi _\sigma ^𝔨^{}.$$ See for an explicit description of the 2-form $`\sigma `$. We can now state a refined version of the Ginzburg-Weinstein theorem . A similar result was obtained by Enriquez-Etingof-Marshall in , for formal Poisson Lie groups. ###### Theorem 4.2 (Ginzburg-Weinstein diffeomorphisms). Let $`K`$ be a compact Lie group with the standard Poisson structure. Then there exists a bisection $`\psi \mathrm{\Gamma }(𝔨^{},K)`$, with $`\psi (0)=1`$, such that the map $$\gamma =e𝒜(\psi ):𝔨^{}K^{}$$ is a Poisson diffeomorphism. The bisection $`\psi `$ can be chosen to be $`T`$-equivariant and to take values in the semi-simple part $`K^{ss}`$. ###### Proof. By Proposition 3.10, there exists a bisection $`\psi \mathrm{\Gamma }(𝔨^{},K),\psi (0)=1`$ with $`\sigma _\psi =\sigma `$. For any such bisection $`𝒜(\psi )_{}\pi ^𝔨^{}=\pi _\sigma ^𝔨^{}=(e^1)_{}\pi ^K^{}`$. Since $`\sigma `$ is $`T`$-invariant, one can take $`\psi `$ to be $`T`$-equivariant. The map $`\psi `$ lifts to a unique map $`\widehat{\psi }\mathrm{\Gamma }(M,\widehat{K}),\widehat{\psi }(0)=1`$ with values in the finite cover $`\widehat{K}=K^{ss}\times Z(K)`$ of $`K`$. Replacing $`\psi `$ with the $`K^{ss}`$-component of $`\widehat{\psi }`$, we arrange that $`\psi `$ takes values in $`K^{ss}`$. ∎ ###### Definition 4.3. A bisection $`\psi \mathrm{\Gamma }(𝔨^{},K)`$ will be called a Ginzburg-Weinstein twist if it has the properties $`\psi (0)=1`$ and $`\sigma _\psi =\sigma `$. By Proposition 3.3(b), Ginzburg-Weinstein twists are unique up to a Lagrangian bisection. Ginzburg-Weinstein twists can be used to turn ordinary $`𝔨^{}`$-valued moment maps into $`K^{}`$-valued moment maps, and vice versa. However, the change of the moment map produces a twisted action. ###### Definition 4.4. Given an $`K\mathrm{Diff}(M)`$ on a manifold $`M`$, and a bisection $`\psi \mathrm{\Gamma }(M,K)`$, the $`\psi `$-twisted action of $`K`$ on $`M`$ is defined as follows, (29) $$𝒜^\psi :K\mathrm{Diff}(M),𝒜^\psi (k)=𝒜(\psi )𝒜(k)𝒜(\psi )^1.$$ ###### Proposition 4.5. Suppose $`\psi \mathrm{\Gamma }(𝔨^{},K)`$ is a Ginzburg-Weinstein twist, and let $`\gamma =e𝒜(\psi )`$. Let $`(M,\pi )`$ be a Poisson manifold, and $`\mathrm{\Phi }:M𝔨^{}`$ and $`\mathrm{\Psi }:MK^{}`$ two Poisson maps related by $`\mathrm{\Psi }=\gamma \mathrm{\Phi }`$. Then $`\mathrm{\Phi }`$ is the moment map for a Hamiltonian $`K`$-action $`𝒜`$ if and only if $`\mathrm{\Psi }`$ is the moment map for a Poisson Lie group $`K`$-action $`𝒜^{}`$. The two actions are related as follows, (30) $$𝒜^{}=𝒜^{\mathrm{\Phi }^{}\psi ^1},𝒜=(𝒜^{})^{(e^1\mathrm{\Psi })^{}\psi }.$$ ###### Proof. Suppose $`\mathrm{\Phi }`$ generates a $`K`$-action $`𝒜`$. We will show that $`\mathrm{\Psi }`$ generates the action $`𝒜^{}=𝒜^{\psi ^1\mathrm{\Phi }}`$. By , the map $$e\mathrm{\Phi }:MK^{}$$ is the moment map for a Poisson-Lie group action of $`(K,\pi ^K)`$ on $`M`$, where $`M`$ is equipped with the gauge transformed Poisson structure $`\pi _{\mathrm{\Phi }^{}\sigma }`$. Since $`\sigma =\sigma _\psi `$, the diffeomorphism $`𝒜(\mathrm{\Phi }^{}\psi ^1)`$ takes the gauge transformed Poisson structure $`\pi _{\mathrm{\Phi }^{}\sigma }`$ structure back to $`\pi `$. Furthermore, $`𝒜(\mathrm{\Phi }^{}\psi ^1)`$ intertwines $`𝒜`$ with the twisted action $`𝒜^{}=𝒜^{\psi ^1\mathrm{\Phi }}`$, and takes $`e\mathrm{\Phi }`$ to $$(e\mathrm{\Phi })𝒜(\mathrm{\Phi }^{}\psi )=e𝒜(\psi )\mathrm{\Phi }=\mathrm{\Psi }.$$ It follows that $`\mathrm{\Psi }`$ is a moment map for the twisted action $`𝒜^{}=𝒜^{\mathrm{\Phi }^{}\psi ^1}`$ on $`(M,\pi )`$. Conversely, assume that $`\mathrm{\Psi }`$ generates an actions $`𝒜^{}`$. Then $`e^1\mathrm{\Psi }:M𝔨^{}`$ is a moment map for a Hamiltonian action on $`(M,\pi _{(e^1\mathrm{\Psi })^{}\sigma })`$. Applying $`𝒜((e^1\mathrm{\Psi })^{}\psi `$ to restore the Poisson structure $`\pi `$, and arguing as above, we see that $`\mathrm{\Phi }`$ is a moment map for the action $`𝒜=(𝒜^{})^{(e^1\mathrm{\Psi })^{}\psi }`$ on $`(M,\pi )`$. (Alternatively, one can also use Lemma 4.9 below to argue that the two formulas (12) are equivalent.) ∎ ### 4.3. Functorial properties of Ginzburg-Weinstein maps A homomorphism of Poisson Lie groups $`K_1,K_2`$ is a Lie group homomorphism $$𝒯:K_1K_2$$ which is also a Poisson map. On the infinitesimal level, a homomorphism of Poisson Lie groups defines a homomorphism of Lie bialgebras, $`\tau :𝔨_1𝔨_2`$. That is, the dual map $`\tau ^{}:𝔨_2^{}𝔨_1^{}`$ is a Lie algebra homomorphism, and in particular exponentiates to a dual Poisson Lie group homomorphism $`𝒯^{}:K_2^{}K_1^{}`$. Given a Poisson Lie group action $`𝒜`$ of $`K_2`$ on a Poisson manifold $`M`$, with moment map $`\mathrm{\Psi }:MK_2^{}`$, the composition $`𝒯^{}\mathrm{\Psi }`$ is a moment map for the $`K_1`$-action $`𝒜𝒯`$. For any compact Lie group $`K`$ with the standard Poisson structure, the maximal torus $`T`$ with the zero Poisson structure is a Poisson-Lie subgroup. That is, the inclusion $`𝒯:TK`$ is a Poisson-Lie group homomorphism. ###### Lemma 4.6. Suppose $`\psi \mathrm{\Gamma }(𝔨^{},K)`$ is a $`T`$-equivariant Ginzburg-Weinstein twist, and let $`\gamma =e𝒜(\psi )`$. Then the following diagram commutes: $$\begin{array}{ccc}𝔨^{}& \underset{\tau ^{}}{}& 𝔱^{}\\ \gamma & & & & \\ K^{}& \underset{𝒯^{}}{}& T^{}\end{array}$$ ###### Proof. Let $`𝒯:TK`$ denote the inclusion, and consider the Poisson map $`\mathrm{{\rm Y}}:𝔨^{}𝔱^{}`$ given as the composition of the Poisson maps $`\gamma :𝔨^{}K^{}`$ and $`𝒯^{}:K^{}T^{}𝔱^{}`$. Proposition 4.5 shows that $`\gamma `$ is a moment map for the twisted $`K`$-action $`𝒜^{\psi ^1}`$, and hence $`\mathrm{{\rm Y}}`$ is a moment map for the twisted $`T`$-action, $`𝒜^{\psi ^1}𝒯`$. Since $`\psi `$ is $`T`$-equivariant, the twisted and untwisted $`T`$-actions coincide. Hence, $`\mathrm{{\rm Y}}`$ and $`\tau ^{}`$ are moment maps for the same $`T`$-action. It follows that their difference is a $`K`$-invariant function $`𝔨^{}𝔱^{}`$. It hence suffices to show that $`\mathrm{{\rm Y}}`$ and $`\tau ^{}`$ coincide on $`𝔱^{}=(𝔨^{})^T𝔨^{}`$ (fixed point set for the coadjoint action of $`T`$ on $`𝔨^{}`$). That is, we have to show that $`\mathrm{{\rm Y}}`$ restricts to the identity map of $`𝔱^{}`$. Since $`\psi `$ is $`T`$-equivariant, it takes $`𝔱^{}=(𝔨^{})^T`$ to $`T=K^T`$ (fixed point set for the conjugation action of $`T`$ on $`K`$). In particular, $`𝒜(\psi )`$ acts trivially on $`𝔱^{}`$, and hence $`\gamma `$ coincides with $`e`$ on $`𝔱^{}𝔨^{}`$. Since $`e:𝔨^{}K^{}`$ restricts to the natural identification $`𝔱^{}T^{}`$, we conclude that $`\mathrm{{\rm Y}}`$ restricts to the identity map of $`𝔱^{}`$. ∎ ###### Theorem 4.7 (Compatible Ginzburg-Weinstein maps). Let $`K_1,K_2`$ be compact Poisson Lie groups with the standard Poisson structure, and $`𝒯:K_1K_2`$ a Poisson Lie group homomorphism. Given any Ginzburg-Weinstein twist $`\psi _1\mathrm{\Gamma }(𝔨_1^{},K_1)`$, there exists a Ginzburg-Weinstein twist $`\psi _2\mathrm{\Gamma }(𝔨_2^{},K_2)`$, for which the diagram (31) $$\begin{array}{ccc}𝔨_2^{}& \stackrel{\tau ^{}}{}& 𝔨_1^{}\\ \gamma _2& & \gamma _1& & \\ K_2^{}& \underset{𝒯^{}}{}& K_1^{}\end{array}$$ with $`\gamma _i=e_i𝒜_i(\psi _i)`$, commutes. Here $`𝒜_i`$ denotes the coadjoint action of $`K_i`$ on $`𝔨_i^{}`$. One can arrange that $`\psi _2`$ takes values in the semi-simple part $`K_2^{ss}`$. ###### Proof. We may assume, passing to a finite cover of $`K_1`$ if necessary, that the semi-simple part $`K_1^{ss}`$ is simply connected. We begin by choosing an arbitrary $`T_2`$-equivariant Ginzburg-Weinstein twist $`\psi _2`$. We will show how to modify $`\psi _2`$ (possibly destroying the $`T_2`$-equivariance), in such a way that the above diagram commutes. The idea is to apply Proposition 3.11 to the Poisson map $$\mathrm{{\rm Y}}=\gamma _1^1𝒯^{}\gamma _2:𝔨_2^{}𝔨_1^{}.$$ To apply this Proposition, we have to verify that the $`𝔷(𝔨_1)^{}`$-component of $`\mathrm{{\rm Y}}`$ is given by a linear map. In fact, we will show that the $`𝔷(𝔨_1)^{}`$-components of $`\mathrm{{\rm Y}}`$ and $`\tau ^{}`$ are equal. Since $`\psi _2`$ is $`T_2`$-equivariant, Lemma 4.6 shows that $`\gamma _2`$ restricts to the natural identification $`𝔱_2^{}T_2^{}`$. Similarly, $`\gamma _1`$ restricts to the natural identification $`𝔷(𝔨_1)^{}Z(K_1)^{}=𝔷(𝔨_1)^{}`$, since this is true for $`e_2`$, and since the action of $`K_1`$ (hence of $`𝒜(\psi _1)`$) on $`𝔷(𝔨_1)^{}`$ is trivial. Since the Poisson bivector of $`K_2`$ vanishes exactly along $`T_2`$, the map $`𝒯`$ must take $`Z(K_1)T_1`$ into $`T_2`$. Hence, the diagram $$\begin{array}{ccccc}𝔨_2^{}& & 𝔱_2^{}& & 𝔷(𝔨_1)^{}\\ \gamma _2& & \gamma _2& & \gamma _1& & \\ K_2^{}& & T_2^{}& & Z(K_1)^{}\end{array}$$ commutes, proving the claim. It follows in particular that the $`𝔷(𝔨_1)^{}`$-component of $`\mathrm{{\rm Y}}`$ is a moment map for the action of $`Z(K_1)K_1`$ via $`𝒯`$. On the other hand, the $`(𝔨_1^{ss})^{}`$-component is a moment map for some action of $`K_1^{ss}`$, since $`K_1^{ss}`$ is simply connected. Hence, $`\mathrm{{\rm Y}}`$ is the moment map for a $`K_1`$-action. By Proposition 3.11, there exists a Lagrangian bisection $`\varphi \mathrm{\Gamma }_0(𝔨_2^{},K_2),\varphi (0)=1`$, with the property $`\mathrm{{\rm Y}}𝒜(\varphi )=\mathrm{{\rm Y}}_0=\tau ^{}`$. That is, replacing $`\psi `$ with $`\psi ^{}=\psi \varphi `$ the diagram (31) commutes. As in the proof of Theorem 4.2, one can arrange that the new $`\psi `$ takes values in $`K_2^{ss}`$, without changing $`𝒜(\psi )`$. ∎ Let us call two Ginzburg-Weinstein twists $`\psi _1\mathrm{\Gamma }(𝔨_1^{},K_1)`$ and $`\psi _2\mathrm{\Gamma }(𝔨_2^{},K_2)`$ compatible (relative to $`𝒯:K_1K_2`$) if the corresponding Ginzburg-Weinstein diffeomorphism $`\gamma _i=e_i𝒜(\varphi _i)`$ define a commutative diagram (31). The compatibility condition is equivalent to a certain equivariance condition, as the following result shows. ###### Theorem 4.8. Suppose $`\psi _1\mathrm{\Gamma }(𝔨_1^{},K_1)`$ and $`\psi _2\mathrm{\Gamma }(𝔨_2^{},K_2)`$ are compatible Ginzburg-Weinstein twists, and put $$\widehat{\psi }_1=𝒯\psi _1(e_1^1𝒯^{}e_2)\mathrm{\Gamma }(𝔨_2^{},K_2).$$ Then the ’ratio’ $`\widehat{\psi }_1^1\psi _2\mathrm{\Gamma }(𝔨_2^{},K_2)`$ is $`K_1`$-equivariant in the sense that it $``$-commutes with all $`𝒯(k)`$ for all $`kK_1`$. One has the formula, (32) $$(\widehat{\psi }_1^1\psi _2)(\mu )=𝒯\left(\psi _1(\tau ^{}\mu )\right)^1\psi _2(\mu ).$$ ###### Proof. Given arbitrary Ginzburg-Weinstein twists $`\psi _1,\psi _2`$, consider again the moment map $`\mathrm{{\rm Y}}=\gamma _1^1𝒯^{}\gamma _2:𝔨_2^{}𝔨_1^{}`$ as in the proof of Theorem 4.7. Suppose $`(M,\pi )`$ is a Poisson manifold, and $`\mathrm{\Phi }:M𝔨_2^{}`$ is the moment map for a Hamiltonian action $`𝒜:K_2\mathrm{Diff}_\pi (M)`$. Let us compute the $`K_1`$-action generated by $`\mathrm{{\rm Y}}\mathrm{\Phi }`$. Using Proposition 4.5 we have, $$\begin{array}{cccc}& \hfill \mathrm{\Phi }:M𝔨_2^{}& \text{ is a moment map for }& 𝒜\hfill \\ & \hfill \gamma _2\mathrm{\Phi }:MK_2^{}& \text{ is a moment map for }& 𝒜^{\mathrm{\Phi }^{}\psi _2^1}\hfill \\ & \hfill 𝒯^{}\gamma _2\mathrm{\Phi }:MK_1^{}& \text{ is a moment map for }& 𝒜^{\mathrm{\Phi }^{}\psi _2^1}𝒯\hfill \\ & \hfill \gamma _1^1𝒯^{}\gamma _2\mathrm{\Phi }:M𝔨_1^{}& \text{ is a moment map for }& (𝒜^{\mathrm{\Phi }^{}\psi _2^1}𝒯)^{\psi _1e_1^1(𝒯^{}\gamma _2\mathrm{\Phi })}.\hfill \end{array}$$ We may re-write the result as $$\begin{array}{cc}\hfill (𝒜^{\mathrm{\Phi }^{}\psi _2^1}𝒯)^{\psi _1e_1^1𝒯^{}\gamma _2\mathrm{\Phi }}& =(𝒜^{\mathrm{\Phi }^{}\psi _2^1})^{\mathrm{\Phi }^{}(𝒯\psi _1e_1^1𝒯^{}\gamma _2)}𝒯\hfill \\ & =(𝒜^{\mathrm{\Phi }^{}\psi _2^1})^{\mathrm{\Phi }^{}(\widehat{\psi }_1𝒜(\psi _2))}𝒯\hfill \\ & =𝒜^{\mathrm{\Phi }^{}(\psi _2^1\widehat{\psi }_1)}𝒯.\hfill \end{array}$$ In the last line, we have used Lemma 4.9 below to write an iterated twist as a single twist. Assume now that $`\psi _1,\psi _2`$ are compatible. The commutativity of the diagram (31) means that $`\mathrm{{\rm Y}}=\tau ^{}`$. In particular, for any Hamiltonian $`K_2`$-space $`(M,\pi ,\mathrm{\Phi })`$, the twisted $`K_1`$-action $`𝒜^{\mathrm{\Phi }^{}(\psi _2^1\widehat{\psi }_1)}𝒯`$ coincides with the untwisted action $`𝒜𝒯`$. By definition of the twisted action, this is equivalent to (33) $$𝒜(\mathrm{\Phi }^{}(\psi _2^1\widehat{\psi }_1))𝒜(𝒯(k))=𝒜(𝒯(k))𝒜(\mathrm{\Phi }^{}(\psi _2^1\widehat{\psi }_1)),kK_1.$$ Apply this result to $`M=K_2\times 𝔨_2^{}`$, with symplectic structure coming from the identification with $`T^{}K_2`$, and with $`\mathrm{\Phi }:(k,\mu )\mu `$ the moment map for the $`K_2`$-action $`𝒜(k)(h,\mu )=(hk^1,k.\mu )`$. Since the map $`\mathrm{\Gamma }(𝔨_2^{},K_2)\mathrm{Diff}(K_2\times 𝔨_2^{}),\psi 𝒜(\mathrm{\Phi }^{}\psi )`$ is 1-1, the above equation implies $`(\psi _2^1\widehat{\psi }_1)𝒯(k)=𝒯(k)(\psi _2^1\widehat{\psi }_1)`$ as desired. The bisection $`\widehat{\psi }_1^1\psi _2`$ may be re-written, using $$\widehat{\psi }_1^1=𝒯\psi _1^1(e_1^1𝒯^{}e_2)=𝒯\psi _1^1𝒜(\psi _1)\tau ^{}𝒜(\psi _2)^1$$ because of the commutativity of the diagram (31). Thus, $$\widehat{\psi }_1^1(𝒜(\psi _2)(\mu ))=𝒯\left(\psi _1(\tau ^{}\mu )\right)^1,$$ and (32) follows. ∎ In the proof we used the following Lemma: ###### Lemma 4.9. Suppose $`M`$ is a $`K`$-manifold with action $`𝒜`$. Let $`\psi \mathrm{\Gamma }_𝒜(M,K)`$ be a bisection relative to the action $`𝒜`$, and $`\varphi \mathrm{\Gamma }_{𝒜^\psi }(M,K)`$ a bisection relative to the twisted action $`𝒜^\psi `$. Then the iterated twist $`(𝒜^\psi )^\varphi `$ can be written as a single twist, $$(𝒜^\psi )^\varphi =𝒜^{\psi 𝒜(\psi )^{}\varphi }.$$ ###### Proof. By definition, $`𝒜^\psi (\varphi )(x)=𝒜(\psi )𝒜(\varphi (x))𝒜(\psi ^1)(x)`$. Hence $$𝒜^\psi (\varphi )=𝒜(\psi 𝒜(\psi )^{}\varphi \psi ^1).$$ Using this formula we calculate, for all $`kK`$, $$\begin{array}{cc}\hfill (𝒜^\psi )^\varphi (k)& =𝒜^\psi (\varphi )𝒜^\psi (k)𝒜^\psi (\varphi )^1\hfill \\ & =𝒜(\psi 𝒜(\psi )^{}\varphi )𝒜(k)𝒜(\psi 𝒜(\psi )^{}\varphi )^1\hfill \\ & =𝒜^{\psi 𝒜(\psi )^{}\varphi }(k).\hfill \end{array}$$ ### 4.4. Anti-Poisson involutions An anti-Poisson involution of a Poisson manifold $`(M,\pi )`$ is an involutive diffeomorphism $`s\mathrm{Diff}(M)`$ reversing the Poisson structure, $`s_{}\pi =\pi `$. An anti-Poisson involution of a Poisson Lie group $`(K,\pi ^K)`$ is an anti-Poisson involution $`s_K`$ of the underlying Poisson manifold which is also an automorphism of the group $`K`$. In this case, $`s_K`$ canonically induces an anti-Poisson involution of the dual Poisson Lie group $`K^{}`$. Suppose $`K`$ is a compact Lie group with standard Poisson structure. Then any anti-linear involution of the Lie algebra $`𝔤=𝔨^{}`$ preserving the Iwasawa decomposition and the bilinear form $`2\mathrm{Im}B^{}`$ defines an anti-Poisson involution $`s_K`$ of $`K`$. Let $`s_𝔨^{}`$ be the induced involution of $`𝔨^{}`$. ###### Lemma 4.10. There exists a Ginzburg-Weinstein twist $`\psi \mathrm{\Gamma }(𝔨^{},K)`$ which, in addition to the Properties from Theorem 4.2, satisfies the equivariance property $$\psi s_𝔨^{}=s_K\psi .$$ ###### Proof. The Ginzburg-Weinstein twist constructed in the proof of Theorem 4.2 has the required equivariance property under involutions. Indeed, the forms $`\sigma _t,a_t`$ on $`𝔨^{}`$, hence also the Moser 1-form $`b_t`$, change sign under $`s_𝔨^{}`$ (see ). It follows that the Moser vector field $`v_t`$ is $`s_𝔨^{}`$-invariant, while the function $`\beta _t`$ is equivariant, $`\beta _ts_𝔨^{}=s_𝔨\beta _t`$. ∎ Example. If $`K=\mathrm{U}(n)`$, the complex conjugation operation $`s_K(A)=\overline{A}`$ is an anti-Poisson involution. The involution $`s_𝔨^{}`$ is complex conjugation on $`𝔨^{}\mathrm{Herm}(n)`$, and $`s_K^{}`$ is complex conjugation on upper triangular matrices or, equivalently, on the space $`P=\mathrm{Herm}^+(n)`$ of positive definite matrices. Compatibility of a Ginzburg-Weinstein twist $`\psi `$ with these involutions just means $`\psi (\overline{A})=\overline{\psi (A)}`$. In particular, $`\psi `$ restricts to a bisection $`\mathrm{Sym}(n)\mathrm{SO}(n)`$. The functoriality properties of Ginzburg-Weinstein maps generalize in the obvious way to the presence of such involutions. Thus, suppose $`K_i,i=1,\mathrm{},n`$ are compact Poisson Lie groups with standard Poisson structure, and $`s_{K_i}`$ are anti-Poisson involutions of $`K_i`$ of the type discussed above. Assume $`𝒯_i:K_iK_{i+1},i=1\mathrm{},n1`$ are Poisson Lie group homomorphisms with $$𝒯_is_{K_i}=s_{K_{i+1}}𝒯.$$ Then the Ginzburg-Weinstein twists $`\psi _{i,t}\mathrm{\Gamma }(𝔨_i^{},K_i)`$ constructed in Theorem 4.7 can be arranged to satisfy, $$\psi _{i,t}s_{𝔨_i^{}}=s_{K_i}\psi _{i,t}.$$ Indeed, the maps obtained in the proof of Theorem 4.7 automatically have this property, since all constructions are compatible with the involutions. It follows that all maps in the commutative diagram (31) intertwine the various involutions. In particular, one obtains a commutative diagram for the fixed point sets of the involutions. ## 5. Gelfand-Zeitlin systems ### 5.1. Thimm actions The following construction of torus actions from non-Abelian group actions appeared in Thimm’s work on completely integrable systems, and was later clarified by Guillemin-Sternberg in . We will present the Thimm actions using the terminology of bisections. Let $`K`$ be a compact Lie group, with maximal torus $`T`$, and let $`𝔨_{\mathrm{reg}}^{}𝔨^{}`$ be the subset of regular elements, that is, elements whose stabilizer is conjugate to $`T`$. Pick a fundamental Weyl chamber $`𝔱_+^{}𝔱^{}`$. Then $`𝔨_{\mathrm{reg}}^{}=K/T\times \mathrm{int}(𝔱_+^{})`$ as $`K`$-manifolds. Restriction of equivariant bisections over $`𝔨_{\mathrm{reg}}^{}`$ to $`\mathrm{int}(𝔱_+^{})`$ defines a group isomorphism, (34) $$\mathrm{\Gamma }(𝔨_{\mathrm{reg}}^{},K)^K\stackrel{}{}\mathrm{\Gamma }(\mathrm{int}(𝔱_+^{}),T).$$ ###### Lemma 5.1. The isomorphism (34) identifies the subgroups of Lagrangian bisections: $`\mathrm{\Gamma }_0(𝔨_{\mathrm{reg}}^{},K)^K\mathrm{\Gamma }_0(\mathrm{int}(𝔱_+^{}),T)`$. ###### Proof. If $`\psi \mathrm{\Gamma }(𝔨_{\mathrm{reg}}^{},K)^K`$ is Lagrangian, then clearly so is its restriction to $`\mathrm{int}(𝔱_+^{})`$. For the converse, suppose $`\psi `$ restricts to a Lagrangian bisection over $`\mathrm{int}(𝔱_+^{})`$. For any $`\xi 𝔨`$ we have $`\iota (\xi )\mu ,(\psi ^1)^{}\theta ^L=\mu ,\xi \mathrm{Ad}_{\psi (\mu )}(\xi )=0`$, since $`\psi (\mu )K_\mu `$. Hence also $$\iota (\xi )\text{d}\mu ,(\psi ^1)^{}\theta ^L=L(\xi )\mu ,(\psi ^1)^{}\theta ^L\text{d}\iota (\xi )\mu ,(\psi ^1)^{}\theta ^L=0.$$ Since on the other hand the pull-back of $`\text{d}\mu ,(\psi ^1)^{}\theta ^L`$ to $`\mathrm{int}(𝔱_+^{})𝔨_{\mathrm{reg}}^{}`$ is zero, this shows $`\text{d}\mu ,(\psi ^1)^{}\theta ^L=0`$. Thus $`\psi `$ is Lagrangian. ∎ Define a group homomorphism (35) $$\chi :T\mathrm{\Gamma }_0(𝔨_{\mathrm{reg}}^{},K)^K$$ by composing the map inverse to (34) with the inclusion $`T\mathrm{\Gamma }_0(\mathrm{int}(𝔱_+^{}),T)`$ as constant bisections. That is, $`\chi (t):𝔨_{\mathrm{reg}}^{}K`$ is the unique $`K`$-equivariant map with $`\chi (t)(\mu )=t`$ for $`\mu \mathrm{int}(𝔱_+^{})`$. Recall that by Lemma 3.7, $`\mathrm{\Gamma }(𝔨_{\mathrm{reg}}^{},K)^K`$ is the center of $`\mathrm{\Gamma }(𝔨_{\mathrm{reg}}^{},K)`$, and that its action on $`𝔨_{\mathrm{reg}}^{}`$ is trivial. In particular, $`\chi (t)`$ acts trivially on $`𝔨_{\mathrm{reg}}^{}`$. Non-trivial actions are obtained by pulling $`\chi (t)`$ back under an equivariant map, $`\mathrm{\Phi }:M𝔨^{}`$. Thus let $`M_0=\mathrm{\Phi }^1(𝔨_{\mathrm{reg}}^{})M`$, and $`\chi _M(t)=\mathrm{\Phi }^{}\chi (t)\mathrm{\Gamma }(M_0,K)^K`$. We define the Thimm action of $`tT`$ by $$tx=𝒜(\chi _M(t))(x),xM_0$$ By construction, the Thimm action commutes with the $`K`$-action, and the map $`\mathrm{\Phi }`$ is Thimm-invariant: $$\mathrm{\Phi }(tx)=t\mathrm{\Phi }(x)=\mathrm{\Phi }(x).$$ From now on, we will write $`\chi (t)(\mu )\chi (t;\mu )`$ and similarly for $`\chi _M`$. ###### Lemma 5.2. If $`\psi \mathrm{\Gamma }(M_0,K)`$ is constant along the fibers of $`\mathrm{\Phi }`$, then $`\psi `$ commutes (under $``$) with all $`\chi _M(t)`$, and $`(\psi \chi _M(t))(x)=\psi (x)\chi _M(t;x)`$. ###### Proof. Since $`𝒜(\chi _M(t))`$ preserves the fibers of $`\mathrm{\Phi }`$, the bisection $`\psi `$ satisfies $`𝒜(\chi _M(t))^{}\psi =\psi `$. Hence, Lemma 3.1 applies. ∎ Thimm actions are naturally associated with Hamiltonian group actions. ###### Lemma 5.3 (Guillemin-Sternberg ). Suppose $`(M,\pi )`$ is a Hamiltonian $`K`$-manifold, with moment map $`\mathrm{\Phi }:M𝔨^{}`$. Then the Thimm $`T`$-action on $`M_0`$ is Hamiltonian, with moment map $`q\mathrm{\Phi }:M𝔱^{}`$. Here $`q:𝔨^{}𝔱_+^{}𝔱^{}`$ is the unique $`K`$-invariant map with $`q(\mu )=\mu `$ for $`\mu 𝔱_+^{}`$. Suppose now that (36) $$K_1\stackrel{𝒯_1}{}K_2\stackrel{𝒯_2}{}\mathrm{}K_n$$ is a sequence of compact Lie groups and homomorphisms, with differentials $`\tau _i:𝔨_i𝔨_{i+1}`$. For $`i<j`$ we will write $`𝒯_i^j=𝒯_{j1}\mathrm{}𝒯_i:K_iK_j`$, with differential $`\tau _i^j:𝔨_i𝔨_j`$. Take the maximal tori $`T_iK_i`$ and positive Weyl chambers $`𝔱_{i,+}^{}`$ to be compatible, in the sense that for all $`i<n`$, $$𝒯_i(T_i)T_{i+1},\tau _i^{}(𝔱_{i+1,+}^{})𝔱_{i,+}^{}.$$ Let $`M`$ be a $`K_n`$-manifold, and $`\mathrm{\Phi }_n:M𝔨_n^{}`$ an equivariant map. Then each $`K_i`$ acts on $`M`$ via $`𝒯_i^n`$, and we obtain a $`K_i`$-invariant map $`\mathrm{\Phi }_i=(\tau _i^n)^{}\mathrm{\Phi }_n:M𝔨_i^{}`$. Let $$M_0=\underset{i=1}{\overset{n}{}}\mathrm{\Phi }_i^1(𝔨_{i,reg}^{}),$$ and define $`\chi _{i,M}:T_i\mathrm{\Gamma }(M_0,K_n)`$ by $$\chi _{i,M}(t_i)=𝒯_i^n\chi _i(t_i)\mathrm{\Phi }_i,t_iT_i$$ where $`\chi _i(t_i)\mathrm{\Gamma }(𝔨_{i,\mathrm{reg}}^{},K_i)`$. ###### Lemma 5.4. The images of the homomorphisms $`\chi _{i,M}:T_i\mathrm{\Gamma }(M_0,K_n)`$ all commute. Hence, they combine to define a group homomorphism $$\chi _M:T_n\times \mathrm{}\times T_1\mathrm{\Gamma }(M_0,K_n).$$ One has the formula, $$\chi _M(t_n,\mathrm{},t_1;x)=\chi _{1,M}(t_1;x)\mathrm{}\chi _{n,M}(t_n;x).$$ ###### Proof. Let $`t_iT_i,t_jT_j`$ where $`i<j`$. The bisection $`\chi _j(t_j)\mathrm{\Gamma }(𝔨_{j,\mathrm{reg}}^{},K_j)`$ is $`K_j`$-equivariant, while $`𝒯_i^j\chi _i(t_i)(\tau _i^j)^{}`$ is constant along the fibers of $`(\tau _i^j)^{}`$. Hence, Lemma 5.2 shows that the two bisections commute under $``$, and that the product $`(𝒯_i^j\chi _i(t_i)(\tau _i^j)^{})\chi _j(t_j)`$ equals the pointwise product. It follows that $`\chi _{i,M}(t_i)`$ and $`\chi _{j,M}(t_j)`$ commute and that the product $`\chi _{i,M}(t_i)\chi _{j,M}(t_j)`$ equals the pointwise product. ∎ We define the Thimm action of $`t=(t_n,\mathrm{},t_1)T_n\times \mathrm{}\times T_1`$ on $`M_0`$ by $$tx=𝒜(\chi _M(t_n,\mathrm{},t_1))(x).$$ If $`(M,\pi )`$ is a Hamiltonian $`K_n`$-space, with moment map $`\mathrm{\Phi }_n`$, then the Thimm action of $`T_n\times \mathrm{}\times T_1`$ on $`M_0`$ is Hamiltonian, with moment map $$(q_n\mathrm{\Phi }_n,\mathrm{},q_1\mathrm{\Phi }_1):M_0𝔱_n^{}\times \mathrm{}\times 𝔱_1^{}.$$ Here $`q_i:𝔨_i^{}𝔱_{i,+}^{}𝔱_i^{}`$ are the unique $`K_i`$-invariant maps with $`q_i(\mu )=\mu `$ for $`\mu 𝔱_+^{}`$. As a special case, the identity map $`\mathrm{\Phi }:𝔨_n^{}𝔨_n^{}`$ gives rise to a Hamiltonian action of $`T_{n1}\times \mathrm{}\times T_1`$ on $$(𝔨_n^{})_0=\underset{i=1}{\overset{n1}{}}((\tau _i^n)^{})^1(𝔨_{i,reg}^{}).$$ (The torus $`T_n`$ is excluded, since its Thimm action is trivial.) ### 5.2. Thimm actions for Poisson Lie groups Let $`K`$ be a compact Lie group with standard Poisson structure, and $`K_{\mathrm{reg}}^{}K^{}`$ the subset of points whose stabilizer under the dressing action of $`K`$ has maximal rank. Since $`e:𝔨_{\mathrm{reg}}^{}K_{\mathrm{reg}}^{}`$ is a $`K`$-equivariant diffeomorphism, any $`K`$-equivariant map $`\mathrm{\Psi }:MK^{}`$ defines a Thimm $`T`$-action, via the composition $`e^1\mathrm{\Psi }`$. Let $`\psi \mathrm{\Gamma }(𝔨^{},K)`$ be a Ginzburg-Weinstein twist, and $`\gamma =e𝒜(\psi )`$. Parallel to Lemma 5.3 we have: ###### Lemma 5.5. Suppose $`M`$ is a Poisson manifold, and $`\mathrm{\Psi }:MK^{}`$ is a moment map for a Poisson Lie group action $`𝒜:K\mathrm{Diff}(M)`$. Then the Thimm $`T`$-action on $`M_0`$ is Hamiltonian, with moment map $$p\mathrm{\Psi }:M_0𝔱^{}.$$ Here $`p=qe^1:K^{}𝔱^{}`$. If $`\mathrm{\Psi }=\gamma \mathrm{\Phi }`$, where $`\mathrm{\Phi }:M𝔨^{}`$ is a moment map for a Hamiltonian $`K`$-action, then the Thimm actions defined by $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ coincide. ###### Proof. As shown in Proposition 4.5, $`\mathrm{\Phi }`$ is the moment map for the twisted action $`𝒜^{\mathrm{\Phi }^{}\psi }`$ on $`M`$. Since $`𝒜_𝔨^{}(\psi )`$ preserves orbits, $`p=qe^1=q\gamma ^1`$. Thus, $`p\mathrm{\Psi }=q\mathrm{\Phi }`$ where $`\mathrm{\Phi }=\gamma ^1\mathrm{\Psi }`$. Thus, Lemma 5.3 identifies $`p\mathrm{\Psi }`$ as the moment map for the Thimm $`T`$-action corresponding to $`\mathrm{\Phi }`$ (relative to the twisted action $`𝒜^{\mathrm{\Phi }^{}\psi }`$). Since the two $`K`$-actions are conjugate under $`𝒜(\mathrm{\Phi }^{}\psi )`$, the same is true for the two Thimm $`T`$-actions. But since $`\chi _M(t)`$ is $`K`$-equivariant, Lemma 5.2 shows $`\mathrm{\Phi }^{}\psi \chi _M(t)\mathrm{\Phi }^{}\psi ^1=\chi _M(t)`$. Hence, the two Thimm actions coincide. ∎ Suppose (36) is a sequence of homomorphisms of Poisson Lie groups $`K_1,\mathrm{},K_n`$, equipped with the standard Poisson structure. Let $`M`$ be a $`K_n`$-manifold, let $`\mathrm{\Psi }_n:MK_n^{}`$ be an equivariant map, and let $`\mathrm{\Psi }_i:MK_i^{}`$ be the composition of $`\mathrm{\Psi }_n`$ with the map $`(𝒯_i^n)^{}:K_n^{}K_i^{}`$. We then obtain commuting Thimm $`T_i`$-actions on $$M_0=\underset{i=1}{\overset{n}{}}\mathrm{\Psi }_i^1(K_{i,\mathrm{reg}}^{}).$$ If $`(M,\pi )`$ is a Poisson manifold, and $`\mathrm{\Psi }_n`$ is the moment map for a Poisson-Lie group action of $`K_n`$, then the Thimm $`T_n\times \mathrm{}\times T_1`$-action on $`M_0`$ is Hamiltonian, with moment map $$(p_n\mathrm{\Psi }_n,\mathrm{},p_1\mathrm{\Psi }_1):M_0𝔱_n^{}\times \mathrm{}\times 𝔱_1^{}.$$ Here $`p_i=q_ie_i^1`$. In particular, we obtain a Hamiltonian $`T_{n1}\times \mathrm{}\times T_1`$-action on $$(K_n^{})_0=\underset{i=1}{\overset{n1}{}}((𝒯_i^n)^{})^1(K_{i,\mathrm{reg}}^{}).$$ By an inductive application of Theorem 4.7, it is possible to choose Ginzburg-Weinstein twists $`\psi _i\mathrm{\Gamma }(𝔨_i^{},K_i)`$, with $`\psi _i(0)=1`$, which are compatible in the sense the resulting diagram $$\begin{array}{ccccccc}𝔨_n^{}& \stackrel{\tau _{n1}^{}}{}& \mathrm{}& \stackrel{\tau _2^{}}{}& 𝔨_2^{}& \stackrel{\tau _1^{}}{}& 𝔨_1^{}\\ \gamma _n& & & & \gamma _2& & \gamma _1& & \\ K_n^{}& \underset{𝒯_{n1}^{}}{}& \mathrm{}& \underset{𝒯_2^{}}{}& K_2^{}& \underset{𝒯_1^{}}{}& K_1^{}\end{array}$$ with $`\gamma _i=e_i𝒜_i(\psi _i)`$ commutes. ###### Proposition 5.6. For any choice of compatible Ginzburg-Weinstein twists $`\psi _i\mathrm{\Gamma }(𝔨_i^{},K_i)`$, the map $`\gamma _n:𝔨_n^{}K_n^{}`$ intertwines the Thimm $`T_{n1}\times \mathrm{}\times T_1`$-actions on $`(𝔨_n^{})_0`$ and $`(K_n^{})_0`$, as well as their moment maps. The map $`\psi _n`$ has the following equivariance property under the Thimm action of $`t=(t_{n1},\mathrm{},t_1)T_{n1}\times \mathrm{}\times T_1`$, (37) $$\psi _n(t\mu )=\stackrel{~}{\chi }(t;\mu )\psi _n(\mu )\chi (t;\mu )^1.$$ Here $$\chi (t;\mu )=\underset{i=1}{\overset{n1}{}}𝒯_i^n\left(\chi _i(t_i;\mu )\right),\stackrel{~}{\chi }(t;\mu )=\underset{i=1}{\overset{n1}{}}𝒯_i^n\left(\mathrm{Ad}_{\psi _i((\tau _i^n)^{}\mu )}\chi _i(t_i;\mu )\right).$$ ###### Proof. For each $`i<n`$ we obtain commutative diagrams (38) $$\begin{array}{ccccc}𝔨_n^{}& \stackrel{(\tau _i^n)^{}}{}& 𝔨_i^{}& \stackrel{q_i}{}& 𝔱_{i,+}^{}\\ \gamma _n& & \gamma _i& & =& & \\ K_n^{}& \underset{(𝒯_i^n)^{}}{}& K_i^{}& \underset{p_i}{}& 𝔱_{i,+}^{}\end{array}$$ It follows that the map $`\gamma _n`$ intertwines the moment maps for the actions of $`T_{n1}\times \mathrm{}\times T_1`$, as well as the actions themselves. By Theorem 4.8, the commutativity of the diagram (38) implies that the bisection $`\widehat{\psi }_i^1\psi _n\mathrm{\Gamma }(𝔨_n^{},K_n)`$ is $`K_i`$-equivariant, and that $$(\widehat{\psi }_i^1\psi _n)(\mu )=𝒯_i^n\left(\psi _i((\tau _i^n)^{}\mu )\right)^1\psi _n(\mu ).$$ The $`K_i`$-equivariance of the bisection $`\psi =\widehat{\psi }_i^1\psi _n`$ implies the Thimm $`T_i`$-equivariance, (39) $$\psi (t_i\mu )=\mathrm{Ad}_{𝒯_i^n\chi _i(t_i;\mu )}\psi (\mu ).$$ Using $`(\tau _i^n)^{}(t_i\mu )=(\tau _i^n)^{}\mu `$, this yields $$\psi _n(t_i\mu )=𝒯_i^n\left(\mathrm{Ad}_{\psi _i((\tau _i^n)^{}\mu )}\chi _i(t_i;\mu )\right)\psi _n(\mu )𝒯_i^n\left(\chi _i(t_i;\mu )\right)^1,$$ proving (37). ∎ Remarks. 1. Throughout this discussion, we can assume that the functions $`\psi _i`$ take values in the semi-simple part $`K_i^{ss}`$. 2. In the presence of anti-Poisson involutions $`s_{K_i}`$ (of the type discussed in Section 4.4) with $`s_{K_{i+1}}𝒯_i=𝒯_is_{K_i}`$, one can assume that the maps $`\psi _i`$ satisfy $`s_{K_i}\psi _i=\psi _is_{𝔨_i^{}}`$. Thus $`\gamma _n`$ restricts to a diffeomorphism between the fixed point sets of $`s_{𝔨_n^{}}`$ and $`s_{K_n^{}}`$, equivariant for the action of $`T_{n1}^{}\times \mathrm{}\times T_1^{}`$, where $`T_i^{}`$ is the fixed point set of the restriction of $`s_{K_i}`$ to $`T_i`$. ### 5.3. The $`\mathrm{U}(n)`$ Gelfand-Zeitlin system Consider the sequence (36) for the special case $`K_i=\mathrm{U}(i)`$, with the standard choice of maximal tori $`T_i=T(i)`$, and with $`𝒯_i^j:U(i)\mathrm{U}(j)`$ the inclusions as the upper left corner (extended by $`1`$’s along the diagonal). Identifying $`𝔲(i)^{}\mathrm{Herm}(i)`$ as above, the standard choice of fundamental Weyl chamber consists of diagonal matrices with decreasing diagonal entries. The maps $`(\tau _i^j)^{}:𝔲(j)^{}𝔲(i)^{}`$ translate into the projection of a Hermitian $`j\times j`$-matrix onto the $`i`$th principal submatrix, and are clearly compatible with these choices of $`𝔱_{i,+}^{}`$. As shown by Guillemin-Sternberg , the Thimm $`T_{n1}\times \mathrm{}\times T_1`$-action for the sequence of projections $$𝔲(n)^{}\mathrm{}𝔲(2)^{}𝔲(1)^{}$$ defines a completely integrable system on $`𝔲(n)^{}`$, and coincides with the Gelfand-Zeitlin system described in Section 1. Let $`U(i)`$ carry the standard Poisson-Lie group structure corresponding to these choices of $`T_i,𝔱_{i,+}^{}`$ and the scalar product $`B_i(A^{},A)=\mathrm{tr}(A^{}A)`$. The bracket on $`𝔲(i)^{}`$ corresponds to its identification with upper triangular matrices, with real diagonal entries. The map $$(\tau _i^j)^{}:𝔲(j)^{}𝔲(i)^{}$$ projects an upper triangular matrix onto the upper left $`i\times i`$ block, and is easily checked to preserve Lie brackets. Hence, $`𝒯_i^j`$ are Poisson-Lie group homomorphisms. The identification $$\mathrm{U}(i)^{}\mathrm{Herm}^+(i),$$ takes the dressing action of $`\mathrm{U}(i)`$ to the action by conjugation. The maps $$(𝒯_i^j)^{}:\mathrm{U}(j)^{}\mathrm{U}(i)^{}$$ are again identified with projection to the upper left corner, both under the identification with positive definite matrices, and under the identification with the group upper triangular matrices with positive diagonal. The Thimm $`T(n1)\times \mathrm{}\times T(1)`$-action for the sequence of maps $$\mathrm{U}(n)^{}\mathrm{}\mathrm{U}(2)^{}\mathrm{U}(1)^{}$$ is Flaschka-Ratiu’s nonlinear Gelfand-Zeitlin system. Let $`\psi _i:𝔲(i)^{}\mathrm{SU}(i)`$ be compatible Ginzburg-Weinstein twists, with $`\psi _i(0)=1`$ and $`\psi _i(\overline{A})=\overline{\psi _i(A)}`$, and let $`\gamma _i:𝔲(i)^{}\mathrm{U}(i)^{}`$ be the corresponding Ginzburg-Weinstein diffeomorphisms. Then $`\psi _n:𝔲(n)^{}\mathrm{SU}(n)`$ has the properties (i)-(iii) listed at the end of Section 2. This finally completes the proof of Theorems 1.1, 1.2, and 1.3. Furthermore, from the uniqueness properties of $`\psi _n`$ (Theorem 1.3), and since $`\gamma _n`$ is a Poisson map by construction, Theorem 1.4 now comes for free. Remark. While all the arguments in this paper were carried out in the $`C^{\mathrm{}}`$-category, we could equally well have worked in the $`C^\omega `$-category of real-analytic maps. In particular, the distinguished 2-form $`\sigma \mathrm{\Omega }^2(𝔨^{})`$ from Section 4.2 is real-analytic, by the explicit formula given in . It follows that the distinguished Ginzburg-Weinstein twist $`\psi `$ for $`\mathrm{U}(n)`$ is not only smooth, but is in fact real-analytic. ### 5.4. Other classical groups We conclude with some remarks on Gelfand-Zeitlin systems for the other classical groups. Consider first the special orthogonal groups $`\mathrm{SO}(n)`$, with the standard choice of maximal tori. Guillemin-Sternberg’s construction for the series of inclusions $$\mathrm{SO}(2)\mathrm{SO}(3)\mathrm{}$$ produces a Gelfand-Zeitlin torus action over an open dense subset of each Poisson manifold $`𝔰𝔬(n)^{}`$. (Not to be confused with the real locus of $`𝔲(n)^{}`$, which does not carry a Poisson structure.) A dimension count confirms that this defines a completely integrable system. On the other hand, for the symplectic groups the series of inclusions $$\mathrm{Sp}(1)\mathrm{Sp}(2)\mathrm{}$$ does not yield a completely integrable system, since the Gelfand-Zeitlin torus does not have sufficiently large dimension. (By a more sophisticated construction, Harada was able to obtain additional integrals of motion in this case.) Consider now the standard Poisson structures on the groups $`\mathrm{SO}(n)`$ and $`\mathrm{Sp}(n)`$. Unfortunately, the inclusions $`\mathrm{SO}(i)\mathrm{SO}(i+1)`$ are not Poisson Lie group homomorphisms, essentially due to the fact that the Dynkin diagram of $`\mathrm{SO}(i)`$ is not a subdiagram of that of $`\mathrm{SO}(i+1)`$. However, the inclusions $`\mathrm{SO}(i)\mathrm{SO}(i+2)`$ are Poisson Lie group homomorphisms, and so are the inclusions $`\mathrm{Sp}(i)\mathrm{Sp}(i+1)`$. By the same discussion as for the unitary groups, one obtains Ginzburg-Weinstein diffeomorphisms $`𝔰𝔬(n)^{}\mathrm{SO}(n)^{}`$ (resp. $`𝔰𝔭(n)^{}\mathrm{Sp}(n)^{}`$) intertwining the resulting (partial) Gelfand-Zeitlin systems. However, in contrast to the unitary groups, there is no simple uniqueness statement in these cases.
warning/0506/cond-mat0506055.html	ar5iv	text	# Circuit theory for noise in incoherent normal–superconducting dot structures ## Abstract We consider the current fluctuations in a mesoscopic circuit consisting of nodes connected by arbitrary connectors, in a setup with multiple normal or superconducting terminals. In the limit of weak superconducting proximity effect, simplified equations for the second-order cross-correlators can be derived from the general counting field theory, and the result coincides with the semiclassical principle of minimal correlations. We discuss the derivation of this result in a multi-dot case. Fluctuations of charge current in mesoscopic structures are in general sensitive to the interactions and the fermionic nature of electrons. In multi-terminal setups, the geometry of the circuit is important for the cross-correlations, and in superconducting heterostructures, also the Andreev reflection, the superconducting proximity effect and transmission properties of NS interfaces need to be accounted for. The general theory for the full counting statistics of current fluctuations in multi-terminal structures was outlined in Ref. Yu. V. Nazarov and Bagrets, 2002. The calculation of the second-order correlators using this theory can be simplified, from complicated $`4\times 4`$ matrix equations to a Kirchoff-type system for scalar parameters, using an approach discussed also, for example, in Refs. Samuelsson and Büttiker, 2002a; Nagaev et al., 2002. In the incoherent case, Nagaev and Büttiker (2001); Samuelsson and Büttiker (2002a); Belzig and Samuelsson (2003) the result coincides with the semiclassical principle of minimal correlations. Samuelsson and Büttiker (2002a); Nagaev et al. (2002) In this paper we show the derivation of this result in a multi-dot system, and consider a few special cases. The theory considers a network of normal ($`𝒯_N`$) and superconducting ($`𝒯_S`$) terminals ($`𝒯=𝒯_N𝒯_S`$) and nodes ($`𝒩`$), connected by connectors. Each connector $`(i,j)`$ is described by its transmission eigenvalues $`T_n^{ij}`$, Blanter and Büttiker (2000) and each node $`j`$ is characterized by a Keldysh Green function $`\stackrel{ˇ}{G}_j`$, which is a $`4\times 4`$ matrix in the Keldysh($`\stackrel{ˇ}{}`$) $``$ Nambu($`\widehat{}`$) space. In the quasiclassical approximation, assuming stationary state and isotropicity, these are only functions of energy, $`\stackrel{ˇ}{G}(\epsilon )`$. The statistics of the current in the circuit is connected to the generating function $`S(\{\chi _k\}_{k𝒯})`$ of charge transfer, which can be found by solving transport equations for the Green functions. In the stationary case at zero frequency, the noise correlations $`\stackrel{~}{S}_{kl}`$ between the fluctuations $`\delta I_k=I_kI_k`$ of currents flowing into the terminals $`k,l𝒯`$ relate to it through Yu. V. Nazarov and Bagrets (2002); Belzig (2003) $$\stackrel{~}{S}_{kl}_{\mathrm{}}^{\mathrm{}}\frac{\text{d}t}{2}\{\delta I_k(t),\delta I_l(0)\}=\frac{e^2}{t_0}\frac{^2S}{\chi _k\chi _l}\|_{\{\chi _j\}=0}.$$ (1) Here, $`t_0`$ is the duration of the measurement, and the equality applies provided this is much larger than the correlation time of the fluctuations. The boundary conditions for transport are assumed such that the terminals are in an internal equilibrium, where the Green function has the form $$\stackrel{ˇ}{G}_{\text{eq}}=\left(\begin{array}{cc}\widehat{R}& \widehat{K}\\ \widehat{0}& \widehat{A}\end{array}\right),\begin{array}{cc}\hfill \widehat{R}& =u\widehat{\tau }_3+vi\widehat{\tau }_2,\widehat{K}=\widehat{R}\widehat{h}\widehat{h}\widehat{A},\hfill \\ \hfill \widehat{A}& =\widehat{\tau }_3\widehat{R}^{}\widehat{\tau }_3,\widehat{h}=f_L+\widehat{\tau }_3f_T.\hfill \end{array}$$ (2) Here, $`u=\left\|\epsilon \right\|/\sqrt{\epsilon ^2\mathrm{\Delta }^2}`$, $`v=\mathrm{sgn}(\epsilon )\sqrt{u^21}`$ are the coherence factors, and $`\mathrm{\Delta }`$ is the superconducting pair amplitude. The functions $`f_T(\epsilon )=1f_0(\epsilon )f_0(\epsilon )`$ and $`f_L(\epsilon )=f_0(\epsilon )f_0(\epsilon )`$ are the symmetric and antisymmetric parts of $`f_0(\epsilon )=[e^{(\epsilon eV)/(k_BT)}+1]^1`$, where $`T`$ is the temperature and $`V`$ the potential of the terminal. We assume $`V=0`$ in all S terminals to avoid time-dependent effects. For calculation of the statistics of the current, the counting field theory additionally specifies the rotation $$\stackrel{ˇ}{G}_k(\chi _k)=e^{i\chi _k\stackrel{ˇ}{\tau }_K/2}\stackrel{ˇ}{G}_{k,\text{eq}}e^{i\chi \stackrel{ˇ}{\tau }_K/2},\stackrel{ˇ}{\tau }_K\stackrel{ˇ}{\tau }_1\widehat{\tau }_3,$$ (3) at each terminal $`k`$, which connects the “counting fields” $`\chi _k`$ to the Green functions. In circuit theory, Yu. V. Nazarov (1999) transport is modeled by the conservation of the matrix current at each node $`i`$ $$\underset{j𝒞}{}\stackrel{ˇ}{I}^{ij}=\stackrel{ˇ}{0},\stackrel{ˇ}{I}^{ij}=\frac{2e^2}{\pi \mathrm{}}\underset{n}{}\frac{T_n^{ij}[\stackrel{ˇ}{G}_j,\stackrel{ˇ}{G}_i]}{4+T_n^{ij}\left(\{\stackrel{ˇ}{G}_i,\stackrel{ˇ}{G}_j\}2\right)}.$$ (4) The sum runs over all nodes and terminals ($`𝒞=𝒯𝒩`$): we assume the convention that $`T_n^{ij}=0`$ for $`i=j`$ and disconnected points. This matrix is related to the observable charge and energy currents by $$I^{ij}=\frac{1}{8e}_{\mathrm{}}^{\mathrm{}}d\epsilon \text{Tr}\left[\stackrel{ˇ}{\tau }_K\stackrel{ˇ}{I}^{ij}\right],I_E^{ij}=\frac{1}{8e^2}_{\mathrm{}}^{\mathrm{}}d\epsilon \epsilon \text{Tr}\left[\stackrel{ˇ}{\tau }_1\stackrel{ˇ}{I}^{ij}\right].$$ (5) Their dependency on $`\{\chi _i\}`$, in turn, describes the generating function of charge transfer: Yu. V. Nazarov and Bagrets (2002) $$\mathrm{d}S(\{\chi _l\})=\frac{t_0}{e}\underset{k𝒯}{}\underset{j𝒞}{}I^{jk}(\{\chi _l\})\mathrm{d}(i\chi _k).$$ (6) Determining the Green functions at the nodes from Eqs. (3,4) and finally applying Eqs. (5,6), one can in principle find the distribution of the fluctuations in the current. However, the problem becomes considerably simpler if one is interested only in the second moment of this distribution, i.e., the current noise as given in Eq. (1). We proceed calculating the noise by assuming that the superconducting proximity effect is negligible, so that the anomalous parts ($`\widehat{\tau }_1,\widehat{\tau }_2)`$ of the functions vanish in each node. Samuelsson and Büttiker (2002a); Belzig and Samuelsson (2003) Then, one can expand the Green function at node $`j`$ to the first order in the counting fields $`\{\chi _k\}`$, in the Nambu-diagonal form: Samuelsson and Büttiker (2002a); Nagaev et al. (2002); Houzet and Pistolesi (2004) $$\stackrel{ˇ}{G}_j=\left(\begin{array}{cc}\widehat{\tau }_3& 2\widehat{h}_j\widehat{\tau }_3\\ 0& \widehat{\tau }_3\end{array}\right)+\underset{k𝒯}{}i\chi _k\left(\begin{array}{cc}\widehat{h}_j\widehat{b}_k^j& 4\widehat{c}_k^j\widehat{b}_k^j\\ \widehat{b}_k^j& \widehat{h}_j\widehat{b}_k^j\end{array}\right)+\mathrm{},$$ (7) where $`\widehat{b}_k^j(\epsilon )=\widehat{1}b_k^j(\left\|\epsilon \right\|)`$, $`\widehat{c}_k^j=c_T^{jk}+\widehat{\tau }_3c_L^{jk}`$ and $`\widehat{h}=f_L+\widehat{\tau }_3f_T`$. This satisfies the quasiclassical normalization $`\stackrel{ˇ}{G}_j^2=\stackrel{ˇ}{1}`$ up to second order in $`\{\chi _k\}`$. For the matrix currents, the above corresponds to the expansion $$\stackrel{ˇ}{I}^{ij}=\left(\begin{array}{cc}\widehat{0}& \widehat{I}_0^{ij}\\ \widehat{0}& \widehat{0}\end{array}\right)+\underset{k𝒯}{}i\chi _k\left(\begin{array}{cc}\widehat{0}& \widehat{I}_{c,k}^{ij}\widehat{I}_{b,k}^{ij}\\ \widehat{I}_{b,k}^{ij}& \widehat{0}\end{array}\right)+\mathrm{}+\stackrel{ˇ}{I}_{\text{coh.}}^{ij}$$ (8) of Eq. (4), where $`\widehat{I}_0`$, $`\widehat{I}_{b,k}`$ and $`\widehat{I}_{c,k}`$ have the structure $`\widehat{I}=\widehat{\tau }_3I(\widehat{\tau }_3\epsilon )`$, due to symmetries in the Nambu space. Here, $`\stackrel{ˇ}{I}_{\text{coh.}}(\{\chi _k\})`$ contains the off-diagonal Nambu-elements, present if $`j`$ corresponds to a superconducting terminal. In what follows, we neglect this coherent part of the current, assuming there are additional decoherence-inducing sink terms in Eq. (4). Samuelsson and Büttiker (2002a); Belzig and Samuelsson (2003) This is valid provided that the Thouless energy describing the inverse time of flight through the node or the connector is much less than the characteristic energy scales of the problem, or, if there is a strong pair-breaking effect in the node, e.g., due to magnetic impurities. One can then consider expansion (8) in detail, assuming a node $`i`$ is connected to a node or terminal $`j`$. This yields four independent equations of conservation: $$\underset{j𝒞}{}I_T^{ij}=0,\underset{j𝒞}{}I_L^{ij}=0,\underset{j𝒞}{}I_{b,k}^{ij}=0,\underset{j𝒞}{}I_{c,T,k}^{ij}=0,$$ (9) in which $`I_T`$ corresponds to the spectral charge current, $`I_L`$ to the energy current, and the last two to a “noise” current, with the symmetric part defined as $`I_{c,T,k}(\epsilon )=I_{c,k}(\epsilon )+I_{c,k}(\epsilon )`$. The corresponding antisymmetric current $`I_{c,L,k}^{ij}`$ is not needed, as we concentrate on the noise in the charge current. The spectral currents have the form $$I_T^{ij}=g_{ij}(f_T^jf_T^i),I_{b,k}^{ij}=g_{ij}(b_k^jb_k^i),$$ (10a) $$I_L^{ij}=\{\begin{array}{cc}0\hfill & \text{for }j𝒯_S\text{ and }\left\|\epsilon \right\|<\left\|\mathrm{\Delta }\right\|\text{,}\hfill \\ g_{ij}(f_L^jf_L^i)\hfill & \text{otherwise.}\hfill \end{array}$$ (10b) Thus, no energy current flows to the superconductors for $`\left\|\epsilon \right\|<\left\|\mathrm{\Delta }\right\|`$. The fourth current is $$\frac{1}{4}I_{c,T,k}^{ij}=g_{ij}(c_T^{ik}c_T^{jk})(b_k^ib_k^j)(s_{ij}(\epsilon )+s_{ij}(\epsilon )),$$ (11) but it can be eliminated, see below. The factors $`g_{ij}`$ and $`s_{ij}(\epsilon )`$ appearing in the expansion can be identified as the conductances and spectral noise densities characteristic of the connectors, and their exact form depends on whether the connector lies between two normal points (NN) or between a normal and a superconducting point (NS). The expressions for the NS case are lengthy, so for simplicity we use here only the limits $`\epsilon \mathrm{\Delta }`$ and $`\epsilon \mathrm{\Delta }`$ for superconducting Green’s functions, effectively neglecting the exact form of the superconducting density of states (DOS). In this approximation, for an NS connector at $`\left\|\epsilon \right\|\left\|\mathrm{\Delta }\right\|`$ or an NN connector, $$\begin{array}{cc}\hfill s_{ij}^{NN}(\epsilon )=\frac{1}{4}g_{ij}^{NN}[& 2(f_L^i+f_T^i)^2(f_L^j+f_T^j)^2\hfill \\ & +F_{ij}^{NN}(f_L^i+f_T^if_L^jf_T^j)^2],\hfill \end{array}$$ (12a) $$g_{ij}^{NN}=\frac{e^2}{\pi \mathrm{}}\underset{n}{}T_n^{ij},F_{ij}^{NN}=\frac{e^2}{g_{ij}^{NN}\pi \mathrm{}}\underset{n}{}T_n^{ij}(1T_n^{ij}).$$ (12b) The result for an NS connector at $`\left\|\epsilon \right\|\left\|\mathrm{\Delta }\right\|`$ is $$s_{ij}^{NS}(\epsilon )=\frac{1}{2}g_{ij}^{NS}[1(f_L^i)^2(f_T^i)^2+F_{ij}^{NS}(f_T^i)^2],$$ (13a) $$g_{ij}^{NS}=\frac{e^2}{\pi \mathrm{}}\underset{n}{}\frac{2(T_n^{ij})^2}{(2T_n^{ij})^2},$$ (13b) $$F_{ij}^{NS}=\left(g_{ij}^{NS}\right)^1\frac{e^2}{\pi \mathrm{}}\underset{n}{}\frac{16(T_n^{ij})^2}{(2T_n^{ij})^4}(1T_n^{ij}),$$ (13c) as found through an expansion of Eq. (4). Naturally, the results above agree with expressions for the noise generated between two terminals, with $`F_{ij}`$ being the differential Fano factor. Blanter and Büttiker (2000); de Jong and Beenakker (1994) The above equations are supplied with the boundary conditions $$b_k^l=\delta _{kl},c_k^l=0,f_k(\epsilon )=f_0(\epsilon ,V_k,T_k),$$ (14) where $`k`$ and $`l`$ are indices of terminals. These can be found by comparing expansion (7) to Eq. (3) (for N terminals), and by examining the expression for $`\stackrel{ˇ}{I}`$ (for S terminals). Finally, Eqs. (1,5,6) yield the result $`\stackrel{~}{S}_{kl}`$ $`={\displaystyle \underset{j𝒞}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\epsilon {\displaystyle \frac{1}{8}}I_{c,T,l}^{kj}`$ $`={\displaystyle \underset{(i,j)}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\epsilon (b_k^ib_k^j)(b_l^ib_l^j)s_{ij}(\epsilon ),`$ (15) for the correlations between terminals $`k`$ and $`l`$. In the last step, we eliminated all $`c_{T,k}^i`$ from the set of equations, which transforms the result to a sum over all connectors $`(i,j)`$ in the circuit. The equations above have a simple physical interpretation. The first two of Eqs. (9) describe the conservation of charge (T) and energy (L) currents at each energy interval $`[\epsilon ,\epsilon +\mathrm{d}\epsilon ]`$. With boundary conditions (14,12,13), they yield distribution functions $`f_L^i`$, $`f_T^i`$ of electrons at the nodes. In addition, one needs to solve from Eqs. (9) the variable $`b_k^i`$, which characterizes the coupling between terminal $`k`$ and node $`i`$. It turns out that this quantity is in fact the characteristic potential introduced for semiclassical multiterminal calculations. Büttiker (1993) Knowing $`f`$, the standard two-terminal relations Beenakker (1997) (12,13) give the spectral noise densities in each connector, and Eq. (15) describes how these couple to the terminals. The final result is similar to the semiclassical result in diffusive metals, Sukhorukov and Loss (1998) and coincides with the result in dot systems, see below. The assumption of all nodes being in the normal state resulted in a simple way to handle superconductors in one special case: first, it takes into account that no energy current enters superconductors at $`\left\|\epsilon \right\|<\left\|\mathrm{\Delta }\right\|`$, and second, assumes that other effects due to superconductivity are localized in only one connector, where both the conductivity and the generated noise are modified. Our last approximation of a piecewise constant superconducting DOS simplifies the resulting expressions. We implicitly assumed above that there is no inelastic scattering which would drive the system towards equilibrium. However, following Ref. Nagaev, 1995, a strong relaxation of the distribution function in a node may be modeled by assuming that $`f_j`$ has the form of a Fermi function. In the case of relaxation due to strong electron-electron scattering, the corresponding potential $`V_j`$ and temperature $`T_j`$ can be determined by taking the two first moments, $`\text{d}\epsilon `$ and $`\text{d}\epsilon \epsilon `$ of Eqs. (9): $$\underset{j𝒞}{}g_{ij}(V_iV_j)=0,^1\frac{3e^2}{k_B^2\pi ^2},$$ (16a) $$\underset{j𝒞𝒯_S}{}g_{ij}[T_i^2T_j^2+^1(V_i^2V_j^2)]=0.$$ (16b) These describe the conservation of charge and energy currents. If some of the nodes are in non-equilibrium, one can define the effective voltages and temperatures so that Eqs. (16) still apply for the whole circuit. In addition, one can model relaxation due to strong electron-phonon coupling by forcing $`T_i`$ coincide with the lattice temperature, so that only $`V_i`$ need to be determined. It is illustrative to note that the quantum-mechanical counting-field theory agrees with the well-known principle of minimal correlations, which is often used in semiclassical calculations. Blanter and Büttiker (2000); Samuelsson and Büttiker (2002a) In a typical model, one has the Langevin equations $$\underset{j𝒞}{}I^{ij}=0,I^{ij}=g_{ij}(V_jV_i)+\delta I^{ij},$$ (17) where $`\delta I^{ij}`$ are the microscopic fluctuations of the current, generated in the connector $`(i,j)`$. Eliminating voltages $`V_i`$ at the nodes and assuming they do not fluctuate at the terminals, one finds the result $$\delta I_k=\underset{(i,j)}{}(b_k^ib_k^j)\delta I^{ij},$$ (18) for the fluctuations $`\delta I_k`$ in the current flowing to terminal $`k`$. Assuming $`\delta I^{ij}`$ are independent and evaluating $`\frac{1}{2}\{\delta I_k,\delta I_l\}`$, one finds Eq. (15). This coincides with the prediction from the counting field theory, for an arbitrary circuit, provided it is understood that $`s_{ij}`$ should be evaluated using the (average) distribution functions at the nodes. These may in general be in non-equilibrium, and should be obtained from a kinetic equation. Moreover, in the incoherent limit, the semiclassical result is correct also in the presence of superconducting terminals. Nagaev and Büttiker (2001); Samuelsson and Büttiker (2002a) The above discussion also clearly shows that an attempt to evaluate the higher correlators of noise using the principle of minimal correlations fails, as this corresponds to truncating expansion (7) after the first two terms. The higher-order semiclassical corrections needed to fix this are discussed for example in Ref. Nagaev et al., 2002. Consider now an example setup that consists of $`M1`$ nodes between two terminals “$`0`$” and “$`M`$” (see inset of Fig. 1), and attempt to calculate its differential Fano factor, $`F_{\text{tot}}\frac{\stackrel{~}{S}_{00}}{eI}\|_{I=0}`$ at zero temperature $`T_0=T_M=0`$. For simplicity of resulting expressions, we assume that all connectors are identical, sharing the same distribution $`\{T_n^{ij}\}`$ of the transmission eigenvalues. First, if both terminals are normal, the application of Eqs. (9,10,12,14,15) is analogous to the semiclassical calculation presented in Ref. Oberholzer et al., 2002, and yields the result $$F_{\text{tot,NN}}=\{\begin{array}{cc}\frac{1}{3}+\frac{3F1}{3M^2},\hfill & \text{for no relaxation,}\hfill \\ \frac{F}{M},\hfill & \text{for }e\text{-ph relaxation,}\hfill \\ \frac{\sqrt{3}}{4},\hfill & \text{for }e\text{-}e\text{ relaxation, }M\mathrm{}\text{.}\hfill \end{array}$$ (19) This shows that the limit $`M\mathrm{}`$ corresponds to the diffusive limit, due to the isotropicity of electron momentum assumed at the nodes. If terminal $`M`$ is superconducting, and relaxation is negligible, we need to apply Eqs. (9,10,12,13,14,15). In this example $`b_k^j`$ are then straightforward to find, $`f_L^j=f_L^0`$ and $`f_T^j=f_T^0b_0^j+f_T^Mb_M^j`$ for $`j=1,\mathrm{},M1`$. Summation in (15) then leads to a simple result $$F_{\text{tot}}=\frac{2}{3}+\frac{(F_{NS}\frac{2}{3})R_{NS}^3+\frac{M1}{2}(F_{NN}\frac{1}{3})R_{NN}^3}{[R_{NS}+(M1)R_{NN}]^3}.$$ (20) Here $`R_{NN}`$, $`F_{NN}`$, $`R_{NS}`$ and $`F_{NS}`$ are the resistances and differential Fano factors of the $`NN`$ and $`NS`$ connectors, as given in Eqs. (12) and (13). The result applies also for $`M=1`$, and in fact, for $`M=2`$ it is valid even if both connectors have differing $`\{T_n^{ij}\}`$. In the limit $`M\mathrm{}`$, the Fano factor again tends towards that of a diffusive contact, showing the doubling of the shot noise. de Jong and Beenakker (1994) Similar calculation shows that for strong inelastic e-ph scattering one has $`F_{\text{tot,e-ph}}F_{NN}/M`$ for large $`M`$. For relaxation due to e-e scattering, in turn, Eq. (16b) first gives the temperature profile $`T_j=[T_0^2+^1(V_0^2V_j^2)]^{1/2}`$, where $`V_j=b_0^jV_0`$. From Eqs. (15,12a) one then finds that $`F_{\text{tot,e-e}}=\sqrt{3}/2`$ for $`M\mathrm{}`$, showing again the doubling of the noise. Numerical results for the behavior at smaller $`M`$ are shown in Fig. 1. It is mostly straightforward to solve the current correlations in multiterminal N-S systems, also discussed for example in Refs. Samuelsson and Büttiker, 2002b; Börlin et al., 2002; Nagaev, 2001. For the four-terminal setup shown in the inset of Fig. 2, one obtains <sup>1</sup><sup>1</sup>1 This type of a problem can also be solved exactly, including the proximity effect, see Ref. Börlin et al., 2002. $`\stackrel{~}{S}_{13}`$ $`=c_1e(\left\|V_1\right\|+\left\|V_3\right\|+\left\|V_1+V_3\right\|)c_2e\left\|V_1V_3\right\|`$ (21) $`c_1`$ $`={\displaystyle \frac{r}{16(r+R)^4}}[rR(1+F)+2FR^2+2r^2(1f)]`$ (22) $`c_2`$ $`={\displaystyle \frac{r}{16(r+R)^4}}[r(2r+R+2r^2R^1)+2r^2f`$ (23) $`+(2+rR^1)(2r^2+4rR+R^2)F].`$ Here, $`R=R_{NN}`$, $`r=R_{NS}`$, $`F=F_{NN}`$ and $`f=F_{NS}`$, and the result is valid provided $`\left\|V_i\right\|\left\|\mathrm{\Delta }\right\|`$, $`i=1,3`$, and $`T=0`$. If the connectors are assumed diffusive ($`F=1/3`$, $`f=2/3`$), Eq. (21) agrees with Ref. Nagaev, 2001. One also finds that for $`f>1`$, the cross-correlation (21) can be positive if $`R`$ is small enough, Samuelsson and Büttiker (2002b) contrary to the case in normal-state circuits. For NS contacts with transparency $`\mathrm{\Gamma }`$, this is satisfied for $`0<\mathrm{\Gamma }<2(\sqrt{2}1)`$, as in Refs. Torrès and Martin, 1999; Samuelsson and Büttiker, 2002b. A different example, where all four contacts are identical so that $`R`$ is not small, is shown in Fig. 2. In conclusion, we discuss a simple model for the transmission of noise in multi-dot incoherent normal–superconducting structures, applying the microscopic counting field theory. The formalism produces the principle of minimal correlations, and has strong analogies with the semiclassical theory of noise in diffusive structures. We thank W. Belzig for discussions, and P. Samuelsson and M. Büttiker for pointing out their previous work in Ref. Samuelsson and Büttiker, 2002a. TTH acknowledges the funding by the Academy of Finland.
warning/0506/hep-ph0506185.html	ar5iv	text	# Short Distance Analysis of 𝑫^𝟎-𝑫̄^𝟎 Mixing ## I Introduction Experimental efforts to detect $`D^0\overline{D}^0`$ mixing are longstanding and remain an active area to this day. data1 ; data2 ; data3 The theory of $`D^0\overline{D}^0`$ mixing is relevant both in lending phenomenological guidance to ongoing experimental work and in better understanding the workings of the Standard Model and of various New Physics scenarios Petrov:2003un ; Burdman:2003rs . In this paper, we present new results — the perturbative QCD NLO contributions in the framework of the $`1/m_c`$ expansion for $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ and $`\mathrm{\Delta }M_D`$. The complex of $`D`$-meson phenomena presents a nontrivial theoretical laboratory for studying applicability of heavy quark methods. One can argue that $`m_c\mathrm{\Lambda }_{\mathrm{QCD}}`$ justifies the use of heavy quark methods. However, the scale $`\mu M_D`$ lies in the meson resonance region, so QCD dynamics is clearly present Golowich:1998pz . As such, there is inherently a degree of interest in the numerical aspect of our findings. Our calculation also touches on matters of principle, such as the degree of $`m_q/m_c`$ suppression in $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ and $`\mathrm{\Delta }M_D`$ at NLO order. We begin by reviewing the theoretical context of $`D^0\overline{D}^0`$ mixing. The mixing arises from $`\mathrm{\Delta }C=2`$ interactions that generate off-diagonal terms in the mass matrix for $`D^0`$ and $`\overline{D}^0`$ mesons. The expansion of the off-diagonal terms in the neutral $`D`$ mass matrix to second order in perturbation theory is $$(M\frac{i}{2}\mathrm{\Gamma })_{12}=\frac{1}{2M_D}\overline{D}{}_{}{}^{0}\|_w^{\mathrm{\Delta }C=2}\|D^0+\frac{1}{2M_D}\underset{n}{}\frac{\overline{D}{}_{}{}^{0}\|_w^{\mathrm{\Delta }C=1}\|nn\|_w^{\mathrm{\Delta }C=1}\|D^0}{M_DE_n+iϵ},$$ (1) where $`_w^{\mathrm{\Delta }C=2}`$ is the effective $`\mathrm{\Delta }C=2`$ hamiltonian and $`_w^{\mathrm{\Delta }C=1}`$ is $$_w^{\mathrm{\Delta }C=1}=\frac{G_F}{\sqrt{2}}\underset{q,q^{}}{}V_{cq}^{}V_{uq^{}}\left[C_1(\mu )Q_1+C_2(\mu )Q_2\right].$$ (2) In $`_w^{\mathrm{\Delta }C=1}`$, the flavor sum on $`q,q^{}`$ extends over the $`d,s`$ quarks,<sup>1</sup><sup>1</sup>1In this paper, we work with $`m_u=m_d=0`$. the quantities $`C_{1,2}(\mu )`$ are Wilson coefficients evaluated at energy scale $`\mu `$, and $`Q_{1,2}`$ are the four-quark operators $$Q_1=\left(\overline{q}_ic_j\right)_{VA}\left(\overline{u}_iq_j^{}\right)_{VA}\mathrm{and}Q_2=\left(\overline{q}_ic_i\right)_{VA}\left(\overline{u}_jq_j^{}\right)_{VA}.$$ (3) The first term in Eq. (1) represents $`\mathrm{\Delta }C=2`$ contributions that are local at scale $`\mu M_D`$, so it contributes to the $`M_{12}`$ (but not to the $`\mathrm{\Gamma }_{12}`$) part of the mixing matrix. For example, in the Standard Model this term is generated by the contribution of the $`b`$ quark. It can also receive a potentially large enhancement from new physics. The second term in Eq. (1) comes from a double insertion of $`\mathrm{\Delta }C=1`$ operators in the SM lagrangian, and it contributes to both $`M_{12}`$ and $`\mathrm{\Gamma }_{12}`$. It is dominated by SM contributions even in the presence of new physics. At scale $`\mu M_D`$, the contributions are from the strange and down quarks and these have relatively large CKM factors. By contrast the $`\mathrm{\Delta }C=2`$ term is expected to give a negligible contribution (e.g. in the SM there is the severe CKM suppression $`\|V_{ub}V_{cb}^{}\|^2/\|V_{us}V_{cs}^{}\|^2=𝒪(10^6)`$). Thus, we omit it henceforth. The off-diagonal mass-matrix terms induce mass eigenstates $`D_L`$ and $`D_S`$ which are superpositions of the flavor eigenstates $`D^0`$ and $`\overline{D}^0`$, $$\|D_{L,S}=p\|D^0\pm q\|\overline{D}{}_{}{}^{0},$$ (4) where $`\|p\|^2+\|q\|^2=1`$. In the Standard Model CP violation in $`D`$ mixing is negligible, as is CP violation in $`D`$ decays both in the Standard Model and in most scenarios of new physics. We therefore assume in the rest of this paper that CP is a good symmetry, and adopt the phase convention Donoghue:1992dd $$𝒞𝒫\|D^0=\|\overline{D}^0,$$ (5) Then we have $`p=q`$, and $`\|D_{L,S}`$ become the CP eigenstates $`\|D_\pm `$ with $`𝒞𝒫\|D_\pm =\pm \|D_\pm `$. We then define the mass and width differences $$\mathrm{\Delta }M_\mathrm{D}M_{D_+}M_D_{}\mathrm{and}\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}\mathrm{\Gamma }_{D_+}\mathrm{\Gamma }_D_{}.$$ (6) It is, however, customary to work directly with the dimensionless quantities, $$x_\mathrm{D}\frac{\mathrm{\Delta }M_\mathrm{D}}{\mathrm{\Gamma }_\mathrm{D}},y_\mathrm{D}\frac{\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}}{2\mathrm{\Gamma }_\mathrm{D}},$$ (7) where $`\mathrm{\Gamma }_\mathrm{D}`$ is the average width of the two neutral $`D`$ meson mass eigenstates. The discussion thus far covers relevant background material. We conclude this section by addressing three particularly important additional points: 1. Our calculation adopts an operator product expansion (OPE) Georgi:1992as ; Bigi:2000wn . In the limit $`m_c\mathrm{\Lambda }`$, where $`\mathrm{\Lambda }`$ is some soft QCD scale, the momentum flowing through the light degrees of freedom in the intermediate state is large. As such, an OPE is implemented by expanding the second term in Eq. (1) in series of matrix elements of local operators. For example, one writes for $`\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}`$, $$\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}=2\mathrm{\Gamma }_{12}=\frac{1}{M_D}\mathrm{Im}\overline{D}{}_{}{}^{0}\|i\mathrm{d}^4xT\left\{_w^{\mathrm{\Delta }C=1}(x)_w^{\mathrm{\Delta }C=1}(0)\right\}\|D^0,$$ (8) and expands the time ordered product in Eq. (8) in local operators of increasing dimension (higher dimension operators being suppressed by powers of $`\mathrm{\Lambda }/m_c`$). 2. We calculate $`\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}`$ by making direct use of work available in the literature Beneke:2003az , but not heretofore applied to $`D^0\overline{D}^0`$ mixing. We then compute the mass difference $`\mathrm{\Delta }M_\mathrm{D}`$ from an unsubtracted dispersion relation,<sup>2</sup><sup>2</sup>2The tiny $`b`$-quark contribution, neglected here, would contribute to a subtraction constant. $$\mathrm{\Delta }M_\mathrm{D}(m_c^2)=\frac{1}{2\pi }\mathrm{P}_{s_0}^{\mathrm{}}ds\frac{\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}(s)}{sm_c^2},$$ (9) which follows from the analyticity of $`\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}`$ in the complex $`s`$-plane with a unitarity branch cut along the $`es`$ axis Falk:2004wg . 3. We expand all our expressions for $`x_{\|rmD}`$ and $`y_\mathrm{D}`$ in powers of the ratio $`z=m_s^2/m_c^2`$. ## II Analysis In what follows we compute LO and NLO contributions to $`y`$ and then $`x`$, $$y_\mathrm{D}=y_{\mathrm{LO}}+y_{\mathrm{NLO}}\mathrm{and}x_\mathrm{D}=x_{\mathrm{LO}}+x_{\mathrm{NLO}}.$$ (10) We depict in Fig. 1 how QCD affects the $`D^0`$-to-$`\overline{D}^0`$ mixing amplitude: (a) the limit of no QCD corrections, (b) the LO component in which QCD dresses the interaction vertices, and (c) an example of a NLO correction. The leading contribution to $`\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}`$ in the $`1/m_c`$ expansion to $`D^0\overline{D}^0`$ mixing comes from the dimension-six $`\|\mathrm{\Delta }C\|=2`$ four-quark operators, $`Q`$ $`=`$ $`\overline{u}_\alpha \gamma _\mu P_Lc_\alpha \overline{u}_\beta \gamma _\mu P_Lc_\beta ,Q_S=\overline{u}_\alpha P_Lc_\alpha \overline{u}_\beta P_Lc_\beta ,`$ $`Q^{}`$ $`=`$ $`\overline{u}_\alpha \gamma _\mu P_Lc_\beta \overline{u}_\beta \gamma _\mu P_Lc_\alpha ,O_S^{}=\overline{u}_\alpha P_Lc_\beta \overline{u}_\beta P_Lc_\alpha ,`$ (11) where $`P_L=(1+\gamma _5)/2`$. One can use Fierz identities and equations of motion to eliminate $`Q^{}`$ and $`Q_S^{}`$ in favor of $`Q`$ and $`Q_S`$. The resulting expression for $`\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}`$ is then $$\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{D}=\frac{G_F^2m_c^2}{12\pi M_D}[F(z)\overline{D}{}_{}{}^{0}\|Q(\mu ^{})\|D^0+F_S(z)\overline{D}{}_{}{}^{0}\|Q_S(\mu ^{})\|D^0],$$ (12) Coefficients $`F(z)`$ and $`F_S(z)`$ are defined as $`F(z)`$ $`=`$ $`{\displaystyle \underset{qq^{}}{}}\xi _q\xi _q^{}\left[F_{11}^{qq^{}}(z)C_1^2(\mu )+F_{12}^{qq^{}}(z)C_1(\mu )C_2(\mu )+F_{22}^{qq^{}}(z)C_2^2(\mu )\right],`$ $`F_{ij}^{qq^{}}(z)`$ $`=`$ $`F_{ij}^{(0)qq^{}}(z)+{\displaystyle \frac{\alpha _s(\mu )}{4\pi }}F_{ij}^{(1)qq^{}}(z),`$ (13) and similarly for $`F_S(z)`$. Here $`\xi _qV_{cq}^{}V_{uq}`$ is a CKM factor for the intermediate $`s,d`$ quarks, the $`\{F_{ij}^{(0)qq^{}}(z)\}`$ functions are given in the discussion to follow and the $`\{F_{ij}^{(1)qq^{}}(z)\}`$ are considered later in our NLO analysis. As usual, the $`D^0`$-to-$`\overline{D}^0`$ matrix elements of $`Q`$ and $`Q_S`$ are parameterized in terms of B-factors, $$\overline{D}{}_{}{}^{0}\|Q\|D^0=\frac{8}{3}f_D^2M_D^2B_\mathrm{D}\mathrm{and}\overline{D}{}_{}{}^{0}\|Q_S\|D^0=\frac{5}{3}f_D^2M_D^2\overline{B}_\mathrm{D}^{(S)},$$ (14) where $`\overline{B}_\mathrm{D}^{(S)}B_\mathrm{D}^{(S)}M_D^2/m_c^2`$. There are limits on the precision of $`B_\mathrm{D}`$ and $`\overline{B}_\mathrm{D}^{(S)}`$ because the calculable short distance component most likely gives a negligibly small contribution. The most recent result for the quenched lattice calculation of $`B_\mathrm{D}`$ is reported in Ref. Gupta:1996yt . ### II.1 Leading Order (LO) Contributions At leading order in $`\alpha _s`$, one finds for the $`s\overline{s}`$ intermediate state contributions to $`F(z)`$ and $`F_S(z)`$, $`\begin{array}{c}F_{11}^{(0)ss}(z)=3\sqrt{14z}(1z)\hfill \\ F_{12}^{(0)ss}(z)=2\sqrt{14z}(1z)\hfill \\ F_{22}^{(0)ss}(z)={\displaystyle \frac{1}{2}}(14z)^{3/2}\hfill \end{array}\begin{array}{c}F_{S11}^{(0)ss}(z)=3\sqrt{14z}(1+2z)\hfill \\ F_{S12}^{(0)ss}(z)=2\sqrt{14z}(1+2z)\hfill \\ F_{S22}^{(0)ss}(z)=\sqrt{14z}(1+2z),\hfill \end{array}`$ (21) and for the $`d\overline{s}`$ and $`s\overline{d}`$ contributions, $`\begin{array}{c}F_{11}^{(0)ds}(z)=3(1z)^2\left(1+{\displaystyle \frac{z}{2}}\right)\hfill \\ F_{12}^{(0)ds}(z)=2(1z)^2\left(1+{\displaystyle \frac{z}{2}}\right)\hfill \\ F_{22}^{(0)ds}(z)={\displaystyle \frac{1}{2}}(1z)^3\hfill \end{array}\begin{array}{c}F_{S11}^{(0)ds}(z)=3(1z)^2\left(1+2z\right)\hfill \\ F_{S12}^{(0)ds}(z)=2(1z)^2\left(1+2z\right)\hfill \\ F_{S22}^{(0)ds}(z)=(1z)^2\left(1+2z\right).\hfill \end{array}`$ (28) In addition, we have $`F_{ij}^{(0)dd}(z)=F_{ij}^{(0)ss}(0)`$. Insertion of Eqs. (II),(21),(28) into Eq. (12) results in the following expression for the leading $`𝒪(z^3)`$ contribution, $`y_{\mathrm{LO}}^{(z^3)}={\displaystyle \frac{G_F^2m_c^2f_D^2M_D}{3\pi \mathrm{\Gamma }_D}}\xi _s^2z^3\left(C_2^22C_1C_23C_1^2\right)\left[B_\mathrm{D}{\displaystyle \frac{5}{2}}\overline{B}_\mathrm{D}^{(S)}\right],`$ (29) where $`\mathrm{\Gamma }_D1.610^{12}`$ GeV is the experimentally determined $`D^0`$ decay rate. The above expression for $`y_{\mathrm{LO}}^{(z^3)}`$ agrees in the no-QCD limit of $`C_1=0`$ and $`C_2=1`$ with that found in the literature dk . Since we expect $`5\overline{B}_\mathrm{D}^{(S)}/2>B_\mathrm{D}`$, it follows that $`y_{\mathrm{LO}}<0`$. An expression for $`x_{\mathrm{LO}}`$ is recovered by inserting $`\mathrm{\Delta }\mathrm{\Gamma }_{\mathrm{LO}}`$ into the dispersion relation of Eq. (9). One disperses in the variable $`m_c^2`$ so that $`z=m_s^2/m_c^2m_s^2/s`$. The functions $`\{F_{ij}^{(0)}(z)\}`$ of Eqs. (21),(28) are employed above the threshold for each intermediate state. Although the dispersion integral diverges separately for each of the $`s\overline{s}`$, $`d\overline{d}`$, $`d\overline{s}`$, $`s\overline{d}`$ intermediate states, the flavor-summed expression for $`\mathrm{\Delta }M_\mathrm{D}`$ is rendered finite by GIM cancellations. All integrals are first evaluated analytically and the results are then expanded in powers of $`z`$. We find that the leading order in the $`z`$-expansion for $`x_{\mathrm{LO}}`$ occurs at $`𝒪(z^2)`$, $`x_{\mathrm{LO}}^{(z^2)}`$ $`=`$ $`{\displaystyle \frac{G_F^2m_c^2f_D^2M_D}{3\pi ^2\mathrm{\Gamma }_D}}\xi _s^2z^2\left[C_2^2B_\mathrm{D}{\displaystyle \frac{5}{4}}(C_2^22C_1C_23C_1^2)\overline{B}_\mathrm{D}^{(S)}\right].`$ (30) As with $`y_{\mathrm{LO}}`$, we again regain the standard no-QCD result Bigi:2000wn ; dk . Terms occurring at next-to-leading order in the $`z`$-expansion are straightforward to determine, and we find $`y_{\mathrm{LO}}^{(z^4)}`$ $`=`$ $`{\displaystyle \frac{G_F^2m_c^2f_D^2M_D}{3\pi \mathrm{\Gamma }_D}}\xi _s^2z^4\left[B_\mathrm{D}\left(C_2^24C_1C_26C_1^2\right){\displaystyle \frac{15}{4}}\overline{B}_\mathrm{D}^{(S)}\left(C_2^22C_1C_23C_1^2\right)\right]`$ $`x_{\mathrm{LO}}^{(z^3)}`$ $`=`$ $`{\displaystyle \frac{G_F^2m_c^2f_D^2M_D}{3\pi ^2\mathrm{\Gamma }_D}}\xi _s^2z^3[{\displaystyle \frac{1}{2}}B_\mathrm{D}(C_2^2+2C_1C_2+3C_1^2)`$ (32) $`\mathrm{ln}z(B_\mathrm{D}{\displaystyle \frac{25}{12}}\overline{B}_\mathrm{D}^{(S)})(C_2^22C_1C_23C_1^2)].`$ Notice that at order $`x_{\mathrm{LO}}^{(z^3)}`$, there is now dependence also on $`\mathrm{ln}z5`$. However, these contributions are quite small relative to those of Eqs. (29),(30). Numerical evaluations for $`y_{\mathrm{LO}}`$ and $`x_{\mathrm{LO}}`$ appear in Table 1. The initial two columns display the leading $`z`$-dependences first with QCD turned off (cf Fig. 1(a)) and then with QCD included via the Wilson coefficients of Eq. (2) (cf Fig. 1(b)). The final column exhibits the exact LO results. The spread of values reflects uncertainties in input parameters (in particular, we have allowed for the range $`B_\mathrm{D}^{(S)}/B_\mathrm{D}=0.81.2`$. The collection of LO results in Table 1 gives rise to several interesting questions, but the most obvious one involves the tiny magnitudes. The main suppression arises from the presence of $`z^3210^7`$ in $`y_{\mathrm{LO}}`$ and $`z^2410^5`$ in $`x_{\mathrm{LO}}`$, even though the expansions for $`F_{ij}(z)`$ and $`F_{Sij}(z)`$ begin at $`𝒪(1)`$. Such $`𝒪(1)`$ contributions would be enormous, but they are in fact cancelled away as are $`𝒪(z)`$ terms. As a result, $`y_{\mathrm{LO}}`$ and $`x_{\mathrm{LO}}`$ are rendered tiny. A numerical by-product of the dependence $`y_{\mathrm{LO}}𝒪(z^3)`$ and $`x_{\mathrm{LO}}𝒪(z^2)`$ is that $`\|x_{\mathrm{LO}}\|\|y_{\mathrm{LO}}\|`$. There is of course a corresponding physics explanation. In the diagrams of Fig. 1, the $`b`$-quark contribution is severly CKM suppressed, so only the light $`s,d`$ quarks propagate on internal legs. Since the mixing amplitude will vanish in the $`m_d=m_s=0`$ limit, the breaking of chiral symmetry and of SU(3) flavor symmetry play crucial roles. Thus, a factor of $`m_s^2`$ comes from an $`SU(3)`$ violating mass insertion on each internal quark line and another from an additional mass insertion on each line to compensate the chirality flip from the first insertion. This mechanism of chiral suppression accounts for the $`z^2`$ dependence of $`x_{\mathrm{LO}}`$. In addition, $`y_{\mathrm{LO}}`$ requires yet another factor of $`m_s^2z`$ to lift the helicity suppression for the decay of a scalar meson into a massless fermion pair. ### II.2 Next-to-Leading Order (NLO) Contributions Any way of reducing the chiral and helicity supression in $`x`$ and $`y`$ should lead to an enhancement. In principle, there are both nonperturbative and perturbative ways to achieve this. One might associate nonperturbative effects with the presence of quark condensates in the QCD vacuum Georgi:1992as ; Bigi:2000wn . These contributions (suppressed by powers of $`1/m_c`$) lead to chirality flip the same way mass insertions do, but have an intrinsic scale of $`\mathrm{\Lambda }1\text{GeV}m_s`$. In the realistic case of not-so-large $`m_c`$, such power suppressions are not always sufficient to ensure the smallness of higher order contributions. Therefore, Eqs. (29) and (30) cannot contribute to leading order in the dual expansion in $`m_s`$ and $`1/m_c`$ if higher order terms in the $`1/m_c`$ expansion contain lower powers of $`z`$ than do $`x_{\mathrm{LO}}`$ and $`y_{\mathrm{LO}}`$. It has already been shown that this is indeed the case Georgi:1992as ; Bigi:2000wn . There are also perturbative QCD corrections to $`x`$ and $`y`$, but these have heretofore not been given serious consideration due to the negligibly small LO values for the $`D^0\overline{D}^0`$ mixing parameters. Even taking into account large scale dependence, the LO result gives a tiny contribution. This has stimulated a shift of attention towards the computation of the long-distance sector with varying degrees of model dependence Donoghue:hh ; Falk:2001hx ; Falk:2004wg . But due to their milder dependence on $`m_s`$, the higher order QCD corrections might be able to give relatively large contributions Petrov:1997ch . This occurs, for example, in the $`cs\gamma `$ short distance amplitude, which receives a huge QCD correction ghmw . In this paper, we consider a specific means by which the helicity supression in $`y`$ can be lifted — a perturbative gluon correction (e.g. as in Fig. 1(c)). Having a perturbative gluon traversing the graph for the correlation function is the same as a well-known effect of lifting of helicity supression which follows from having three particles in the intermediate state instead of two <sup>3</sup><sup>3</sup>3This mechanism leads to the prediction that the rate for the weak radiative decay $`B\gamma e\nu `$ is much larger than the rate of weak leptonic decay $`Be\nu `$. The addition of the ‘intermediate-state’ gluon can lift one power of $`z`$, which characterizes the helicity suppression in $`y`$. Also, the relative lightness of $`m_c`$ implies that higher order perturbative QCD corrections are suppressed by a relatively large factor of $`\alpha _s(m_c)0.4`$. It is therefore expected that the NLO corrections to $`y`$ should dominate the LO result. Moreover, the existence of a dispersion relation implies that $`x`$ might well be enhanced at NLO. In order to systematically include the effects of intermediate-state gluons, a complete calculation of NLO corrections to $`D^0\overline{D}^0`$ mixing is needed. The NLO corrections to lifetime difference $`y`$ can be readily computed. All the relevant NLO contributions to $`F(z)`$ and $`F_S(z)`$ for two massive quarks and one massive, one massless quark can be found by adopting the formulas in Refs. Beneke:2003az (which considered the case of $`B_s\overline{B}_s`$ mixing) to computing $`F_{ij}^{(1)qq^{}}`$ of Eq. (II). That calculation has been performed in the NDR-scheme (dimensional regularization with anti-commuting $`\gamma _5`$ and $`\overline{MS}`$ subtraction). We shall not present explicit formulas for the $`\{F_{ij}^{(1)qq^{}}(z)\}`$ and $`\{F_{S,ij}^{(1)qq^{}}(z)\}`$ functions as they are rather cumbersome. Scale dependent quantities used in our numerical work and evaluated at $`\mu =1.3`$ GeV were: $$m_c=1.3\mathrm{GeV},B_\mathrm{D}=0.82,C_1=0.411,C_2=1.208,\alpha _s=0.406.$$ (33) The value for $`B_\mathrm{D}`$ at scale $`\mu =1.3`$ GeV is obtained by referring the lattice determination at $`\mu =2`$ GeV and employing the scale invariant quantity $`\widehat{B}_\mathrm{D}`$, $$\widehat{B}_\mathrm{D}=B_\mathrm{D}(\mu _0)[\alpha _s(\mu _0)]^{6/25}\left[1+\frac{\alpha _s(\mu _0)}{4\pi }J_4\right],$$ (34) with $`J_41.792`$. Also we allow for a range of the ratio $`\overline{B}_\mathrm{D}^{(S)}/B_\mathrm{D}`$. Using the $`\{F_{ij}^{(1)qq^{}}(z)\}`$ and $`\{F_{S,ij}^{(1)qq^{}}(z)\}`$ functions, we have calculated $`y_{\mathrm{NLO}}`$ exactly and also have expressed it in terms of a power series in $`z`$. The leading term is $`𝒪(z^2)`$, $`y_{\mathrm{NLO}}^{(2)}`$ $`=`$ $`{\displaystyle \frac{G_F^2m_c^2f_D^2M_D}{3\pi \mathrm{\Gamma }_D}}\xi _s^2{\displaystyle \frac{\alpha _s}{4\pi }}z^2(B_\mathrm{D}[({\displaystyle \frac{77}{6}}{\displaystyle \frac{8\pi ^2}{9}})C_2^2+14C_1C_2+8C_1^2]`$ (35) $`{\displaystyle \frac{5}{2}}\overline{B}_\mathrm{D}^{(S)}[({\displaystyle \frac{8\pi ^2}{9}}{\displaystyle \frac{25}{3}})C_2^2+20C_1C_2+32C_1^2]),`$ and the corresponding $`𝒪(z^3)`$ contribution is $`y_{\mathrm{NLO}}^{(3)}={\displaystyle \frac{2G_F^2m_c^2f_D^2M_D}{3\pi \mathrm{\Gamma }_D}}\xi _s^2{\displaystyle \frac{\alpha _s}{4\pi }}z^3`$ $`\times (B_\mathrm{D}[(15+7\mathrm{ln}z)C_2^2({\displaystyle \frac{77}{9}}+{\displaystyle \frac{103}{3}}\mathrm{ln}z)C_1C_2(18+58\mathrm{ln}z)C_1^2]`$ $`{\displaystyle \frac{5}{2}}\overline{B}_\mathrm{D}^{(S)}[({\displaystyle \frac{28}{3}}+6\mathrm{ln}z)C_2^2+({\displaystyle \frac{49}{9}}{\displaystyle \frac{118}{3}}\mathrm{ln}z)C_1C_2({\displaystyle \frac{31}{3}}+58\mathrm{ln}z)C_1^2]).`$ (36) The numerical results, displayed in Table II, reveal that $`y_{\mathrm{NLO}}`$ is almost an order of magnitude larger than $`y_{\mathrm{LO}}`$ and that the subleading term $`y_{\mathrm{NLO}}^{(3)}`$ is smaller than $`y_{\mathrm{NLO}}^{(2)}`$ but not at all negligible. The corresponding expression for $`x_{\mathrm{NLO}}`$ has, as before, been obtained by means of a dispersion relation. We evaluated the dispersion integral numerically to obtain the value presented in Table II. As regards an analytical expression for $`x_{\mathrm{NLO}}`$, the intent was again to by first exactly perform the dispersion integrals and then expand each contribution in a $`z`$ power series. It turned out possible to do this for the $`d\overline{d}`$, $`d\overline{s}`$ and $`s\overline{d}`$ intermediate states, but not for $`s\overline{s}`$. It is, however, nonetheless useful to have an approximate analytic representation for $`x_{\mathrm{NLO}}`$. By exploring a variety of approximation techniques, we found the expected $`𝒪(z^2,z^2\mathrm{ln}z)`$ leading behavior for $`x_{\mathrm{NLO}}`$ but encountered scatter in the $`𝒪(z^2)`$ coefficients, although less so for the $`𝒪(z^2\mathrm{ln}z)`$ coefficients. Upon accepting the latter and fitting the $`𝒪(z^2)`$ coefficients to the numerical evaluations of individual dispersion integrals, we arrived at the ‘effective’ formula: $`x_{\mathrm{NLO}}{\displaystyle \frac{G_F^2m_c^2f_D^2M_D}{3\pi ^2\mathrm{\Gamma }_D}}\xi _s^2{\displaystyle \frac{\alpha _s}{4\pi }}z^2`$ $`\times (B_\mathrm{D}[(11.34.1\mathrm{ln}z)C_2^2+(49.2+15.8\mathrm{ln}z)C_1C_2+(37.9+10.7\mathrm{ln}z)C_1^2]`$ (37) $`{\displaystyle \frac{5}{8}}\overline{B}_\mathrm{D}^S[(37.9+2.2\mathrm{ln}z)C_2^2+(33.+81.8\mathrm{ln}z)C_1C_2+(32.0+125.3\mathrm{ln}z)C_1^2]).`$ This relation, although approximate, is nonetheless useful in understanding the magnitude of the various contributions to $`x_{\mathrm{NLO}}`$. Since the NLO results found for the box contributions are larger than their LO counterparts, we consider here for the sake of completeness the NLO penguin contribution $`y_{\mathrm{NLO}}^{(\mathrm{P})}`$ to the width difference. We have $`y_{\mathrm{NLO}}^{(\mathrm{P})}={\displaystyle \frac{4G_F^2m_c^2f_D^2M_D}{9\pi \mathrm{\Gamma }_D}}\xi _s^2{\displaystyle \frac{\alpha _s}{4\pi }}z^3C_2^2\left(B_\mathrm{D}+5\overline{B}_\mathrm{D}^S\right)+\mathrm{}.`$ (38) The result shown in Table 2 clearly shows the penguin amplitude for $`y_{\mathrm{NLO}}^{(\mathrm{P})}`$ is negligible compared to the box contribution. The mass splitting $`x_{\mathrm{NLO}}^{(\mathrm{P})}`$ is likewise $`𝒪(z^3)`$ and hence negligible. ## III Concluding Comments We have calculated LO and NLO contributions to the leading dimension-six component in the OPE for $`D^0\overline{D}^0`$ mixing. Numerical results appear in Table 1 for LO and Table 2 for NLO. As a partial check of our analysis, we found our results (in cases of overlap) to agree with work carried out previously. Our formulae for $`x`$ and $`y`$ involve not simply expansions in $`1/m_c`$, but rather combined expansions in $`m_s`$ ($`m_d`$ is negligible), $`\alpha _s`$, and $`1/m_c`$. As a technical aside, we performed the calculations at scale $`m_c1.3`$ GeV. The two most noteworthy numerical features found for $`x`$ and $`y`$ are: 1. They are small at LO and even at NLO. This is because $`zm_s^2/m_c^2`$ is small and the leading dependence on $`z`$ is found to be $$y_{\mathrm{LO}}z^3x_{\mathrm{LO}}z^2y_{\mathrm{NLO}}z^2x_{\mathrm{NLO}}z^2.$$ (39) Although contributions from individual intermediate states are not small, CKM factors cancel away the $`𝒪(1)`$ and $`𝒪(z)`$ components. 2. The NLO terms are larger than the LO terms. This requires somewhat more explanation, especially since NLO amplitudes contain the small perturbative QCD factor $`\alpha _s/4\pi `$. As regards the dimensionless width difference $`y_\mathrm{D}`$, the ratio of leading terms in the $`z`$ expansion is $$\frac{y_{\mathrm{NLO}}^{(z^2)}}{y_{\mathrm{LO}}^{(z^3)}}=\frac{\alpha _s}{4\pi }\times \frac{1}{z}\times \frac{W_y^{(\mathrm{NLO})}}{W_y^{(\mathrm{LO})}}0.03\times 169\times (0.73)4.$$ (40) In the above $`W_y^{(\mathrm{NLO})}/W_y^{(\mathrm{LO})}`$ is the ratio of terms containing the Wilson coefficients in Eqs. (29),(35) and for definiteness we have considered the case $`B_\mathrm{D}^{(S)}=0.8B_\mathrm{D}`$. Eq. (40) shows that $`\|y_{\mathrm{NLO}}^{(z^2)}\|`$ exceeds $`\|y_{\mathrm{LO}}^{(z^3)}\|`$ because the extra factor of $`z`$ overwhelms the $`\alpha _s/4\pi `$ suppression. We have already discussed the physics of this – the helicity suppression mechanism which affects any LO $`q\overline{q}`$ intermediate state is removed via the presence of a virtual gluon in the NLO $`q\overline{q}G`$ intermediate state. Also, the difference in sign between $`y_{\mathrm{NLO}}^{(z^2)}`$ and $`y_{\mathrm{LO}}^{(z^3)}`$ arises from the factor $`W_y^{(\mathrm{NLO})}/W_y^{(\mathrm{LO})}`$. Since, to leading order in $`z`$, $`x_{\mathrm{NLO}}`$ and $`x_{\mathrm{LO}}`$ both behave as $`z^2`$, something else must account for the result $`\|x_{\mathrm{NLO}}\|>\|x_{\mathrm{LO}}\|`$. From Eq. (30) and the approximate formula Eq. (37), we have $$\frac{x_{\mathrm{NLO}}}{x_{\mathrm{LO}}^{(z^2)}}\frac{\alpha _s}{4\pi }\times \frac{W_x^{(\mathrm{NLO})}}{W_x^{(\mathrm{LO})}}0.03\times (41.4)1.3,$$ (41) where $`W_x^{(\mathrm{NLO})}`$ and $`W_x^{(\mathrm{LO})}`$ are again the contributions from the Wilson constants and their coefficients. In this case, the suppression in $`\alpha _s/4\pi `$ is overcome by the large size of $`W_x^{(\mathrm{NLO})}/W_x^{(\mathrm{LO})}`$. In particular, the largest contributor to $`W_x^{(\mathrm{NLO})}`$ is from the $`\overline{B}_\mathrm{D}^S`$ term in Eq. (37), roughly equally between log and non-log terms. 3. We conclude that, citing just central values, the net effect of the short distance contributions is $$y_\mathrm{D}=y_{\mathrm{LO}}+y_{\mathrm{NLO}}610^7,x_\mathrm{D}=x_{\mathrm{LO}}+x_{\mathrm{NLO}}610^7.$$ (42) In brief, $`y_\mathrm{D}`$ is given by $`y_{\mathrm{NLO}}`$ to a reasonable approximation but $`x_\mathrm{D}`$ is greatly affected by destructive interference between $`x_{\mathrm{LO}}`$ and $`x_{\mathrm{NLO}}`$. The net effect is to render $`y_\mathrm{D}`$ and $`x_\mathrm{D}`$ of similar magnitudes, at least at this order of analysis. $`D^0\overline{D}^0`$ mixing thus provides a concrete example of a well-defined observable for which NLO perturbative QCD corrections dominate the LO result. Will it follow that the NNLO contributions are larger still? Of course, one cannot know without doing the calculation. We feel, however, it may not necessarily be the case, at least as regards the width difference $`y`$. The three-particle $`q\overline{q}G`$ intermediate states were able to lift the helicity suppression experienced by $`q\overline{q}`$ intermediate states. In passing to the NNLO sector, there is no analogous suppression factor to be lifted. Of course, there is always the possibility that large numerical coefficients can overturn the $`\alpha _s/4\pi `$ counting. The question remains – just how large is $`D^0\overline{D}^0`$ mixing? Evidently, it is still not possible to provide a definitive theoretical answer and experiment will presumably decide the issue. On a relative basis, our ‘short-distance’ numerical results are smaller than most ‘long-distance’ estimates (although the model dependence and uncertainty present in even modern and improved versions of the latter is less significant here). Experimentalists might find it useful to interpret our numerical NLO values as lower bounds to $`y_\mathrm{D}`$ and $`x_\mathrm{D}`$. We conclude by considering our analysis in the context of operator product expansions. As we have seen above, the prediction of $`x`$ and $`y`$ is a result of expanding the correlation function Eq. (8) in terms of three ‘small’ quantities, $`z`$, $`\mathrm{\Lambda }/m_c`$, and $`\alpha _s`$. Since the first quantity is significantly smaller than the other two, the structure of the series is rather different from other (usual) applications of the OPE, e.g. $`B^0\overline{B}^0`$ mixing or $`b`$-hadron lifetimes Gabbiani:2004tp . Working with this combined expansion, we computed the leading contribution originating from matrix elements of dimension-six operators. These matrix elements are commonly parameterized in terms of the two nonperturbative parameters, $`B_D`$ and $`\overline{B}_D`$. The applications of techniques of lattice and QCD sum rule evaluations of these operators can hopefully further improve the precision of our prediction. At higher orders in this expansion one would need to take into account $`𝒪(z^{3/2})`$ corrections (multiplied by about a dozen matrix elements of dimension-nine operators) and $`𝒪(z)`$ corrections (with more than twenty matrix elements of dimension-twelve operators). This would introduce a veritable multitude of unknown parameters whose matrix elements cannot be computed at this time. Simple dimensional analysis Bigi:2000wn suggests magnitudes $`x_Dy_D10^3`$, but order-of-magnitude cancellations or enhancements are possible. However, any effect of higher orders in $`1/m_c`$ or $`\alpha _s(m_c)`$ which could render the result to be proportional to $`z^n`$ in the lowest possible power $`n=1`$ Falk:2001hx would presumably produce a dominant contribution to the prediction of $`x`$ and $`y`$. ###### Acknowledgements. The work of E.G. has been supported in part by the U.S. National Science Foundation under Grant PHY–0244801. A.P. was supported in part by the U.S. National Science Foundation under Grant PHY–0244853, and by the U.S. Department of Energy under Contract DE-FG02-96ER41005. We thank Sandip Pakvasa for conversations at the initial stages of this project.
warning/0506/astro-ph0506295.html	ar5iv	text	# The optical counterpart of XTE J0929–314, the third transient millisecond X-ray pulsar ## 1 Introduction X-ray heating of three regions is generally believed to contribute to optical variability in low mass X-ray binary (LMXB) systems. These are the accretion disc, a bright spot on the outer edge of the accretion disc due to inflowing material and the hemisphere of the companion facing the neutron star. In most LMXBs the reprocessed X-ray optical flux dominates the optical light from the rest of the system (van Paradijs 1983, van Paradijs & McClintock 1995), particularly in the outburst phase. The companion itself may only be evident at a very faint level when the system is in quiescence. More recently it has become apparent that synchrotron emission from matter flowing out of the system via bipolar jets makes a highly variable contribution to radio and IR emission from many different classes of X-ray binaries (Fender 2003). In some cases this emission may extend into the optical region (Hynes et al. 2000). At least one other persistent millisecond X-ray pulsar, SAX J1808.4–3658, is known to have a transient IR excess and radio emission probably due to synchrotron processes (Wang et al. 2001). On 2002 April 30 an X-ray transient was discovered by Remillard et al. (2002) using the All Sky Monitor (ASM) (Levine et al. 1996) on The Rossi X-ray Timing Explorer RXTE satellite. XTE J0929–314, the subject of this paper, was subsequently found to also be a millisecond X-ray pulsar by Remillard, Swank & Strohmayer (2002). Using additional RXTE Proportional Counter Array (PCA) observations Galloway et al. (2002a; 2002b) reported a neutron star spin frequency of 185 Hz, a binary period of 2615-s and an implied companion mass of $`0.008M_{}`$, about 8.5 Jupiter masses. A blue and variable optical couterpart was suggested by Greenhill, Giles & Hill (2002). This identification was supported by spectra obtained by Castro-Tirado et al. (2002) who found a number of emission lines superimposed on a blue continuum which is typical of soft X-ray transients in outburst. A coincident radio source was also reported by Rupen, Dhawan & Mioduszewski (2002). XTE J0929–314 was the third transient millisecond X-ray pulsar to be discovered. The first (SAX J1808.4–3658) has been studied extensively at all wavelengths (Giles, Hill & Greenhill 1999; Wang et al. 2001; Homer et al. 2002; Wachter et al. 2000; in ’t Zand et al. 1998; Wijnands et al. 2001; Markwardt, Miller & Wijnands 2002; Wijnands & van der Klis 1998; Chakrabarty & Morgan 1998; Chakrabarty et al. 2003; Wijnands et al. 2003). The V band flux from SAX J1808.4–3658 decayed from $`16.7518.5`$ mag in 10-d suggesting an e-folding time of 5-6 d. Of the other four known millisecond X-ray pulsars only two have identified visible counterparts. XTE J1814–338 (Strohmeyer et al. 2003) has been identified by Krauss et al. (2003) but no detailed optical or IR observations are available. The more recent IGR J00291+5934 (Galloway et al. 2005) has a detailed optical light curve (Bikmaev et al. 2005). The R band flux from this object decayed from $`17.422.4`$ mag in 30-d giving an e-folding time of $`5.66\pm 0.2`$ d. Study of all these systems is expected to provide important information on the evolutionary path by which a conventional LMXB system might turn into a millisecond radio pulsar. In this paper we describe the optical variability of XTE J0929–314. The measurements were made in the BVRI bands during a period of $``$ 9 weeks following its discovery. ## 2 Observations All the observations described in this paper were made using the 1-m telescope at the University of Tasmania Mt. Canopus Observatory. The CCD camera, its operating software (CICADA), the image reduction and analysis tools (MIDAS and DoPHOT) were identical to that described in Giles et al. (1999). The CCD camera contains an SITe chip which is a thinned back illuminated device providing 512 x 512 pixels with an image scale of $`0.434\mathrm{}`$ pixel<sup>-1</sup>. Cousins standard BVRI filters (Bessell 1990) were used for the observations. The data were calibrated using a sequence of observations of 11 standard stars within the RU149D and PG1047 fields of Landolt (1992) to derive the magnitudes of five local secondary standards close to XTE J0929–314 and within the CCD frame. These local standards are marked as stars 1-5 on the finder chart in Fig. 1 and we tabulate their derived magnitudes in Table 1. The magnitudes for XTE J0929–314 were then obtained using differential photometry relative to these local secondary standards. We did not use the same stars for all colours but used a combination of three of the five as indicated in Table 1. This multi-star process revealed that differential colour corrections are negligible. A few observations were interrupted or terminated early by the arrival of clouds. The complete data set from 2002 May 1 to July 1 \[HJD 2452(396) - 2452(457)\] is detailed in Tables 2 & 3 and forms the subject of this paper. ## 3 Results ### 3.1 Source position To determine the source position we selected the nearest 9 stars to the candidate from the Hubble Guide Star Catalogue 2.2 downloaded from the NASA HEASARC web site. We used the CCD coordinates of these stars and the proposed candidate on the best quality I band image from May 2 (HJD 397) to derive an accurate source position for XTE J0929–314 of R.A. 9h 29m 20$`\stackrel{s}{.}19`$ Dec. $`31\mathrm{°}23\mathrm{}3\stackrel{}{.}2`$ with a relative error of $`\pm 0\stackrel{}{.}1`$ (equinox J2000.0). This is 0$`\stackrel{}{.}7`$ from our initial position estimate in Greenhill et al. 2002 and only $``$0$`\stackrel{}{.}2`$ away from the radio position given by Rupen et al. (2002). Chandra observations confirmed that the X-ray source was $``$ 1$`\stackrel{}{.}25`$ from our optical position (Juett, Galloway & Chakrabarty 2003). We note that this difference is twice their quoted error. In Fig. 1 we provide a finder chart for XTE J0929–314 constructed from an I band CCD image from May 2. The object is the central star in a line of three and for the first few weeks was seen to be of similar brightness to its two neighbours. ### 3.2 The X-ray light curve XTE J0929–314 is one of the faintest transients to be found by the ASM experiment on RXTE and this detection was only possible due to its angular distance from the Galactic Centre region which minimises bright source confusion and consequent positional uncertainties. Our optical candidate was only $`0\stackrel{}{.}5`$ from the X-ray position of Remillard et al. (2002). The transient was discovered near the time of peak X-ray flux, towards the end of April. It was then observed periodically with the PCA experiment on RXTE \[obs id 70096-03-(\-\)\] as part of a proprietary TOO campaign (Galloway et al. 2002b; Juett et al 2003). In the top panel of Fig. 2 we show the ASM intensity history for this transient up to just past the time when regular PCA observations commenced. Note that there is a single early PCA observation at HJD 397.0. The ASM points are plotted as daily averages since the data for individual dwell cycles are too noisy for such a relatively faint source. The 38 PCA values are averages over individual PCA observation id’s which typically last 1000-4000 seconds. The two X-ray light curves in Fig 2. are similar to fig. 1 in Galloway et al. (2002b) except that both vertical axes are now drawn with ’X-ray magnitude’ scales. The final three PCA observations are not plotted since only upper limit detections exist after HJD 445.0. ### 3.3 The optical light curves In Fig. 2 we plot our sets of BVRI band data from Tables 2 & 3. The tight clusters of I data on May 2, 3 & 4 (HJD 397, 398 & 399) are shown in greater detail in Fig. 4. There is evidence of variability in the range 0.05 to 0.1 magnitudes on a timescale of hours during several nights when long runs of I band measurements were made. Greenhill et al (2002) reported variability of up to $`0.5`$ magnitudes on the first night (HJD 396) of our observations. Preliminary analysis suggested that a short duration (timescale $`30`$ minutes) $`0.5`$ magnitude flare occurred in V on this night. Subsequently we became aware of water vapour condensation on the filter during this ”event” and we are now doubtful of the reality of the flaring. In the following section we describe evidence for a low amplitude orbital period modulation seen in the I band flux at about HJD 397. It is of interest to note that no significant variations were seen in X-rays during the RXTE PCA observations (Galloway et al. 2002b). The upper limit to the orbital period modulation of the 2-10 keV X-ray flux was $`<1.1`$ per cent (3$`\sigma `$) (Juett et al. 2003). In Fig. 3 we plot the mean B-V and V-I colour indices for the nights when these three colours were measured. In order to minimise the effects of source variability the colour indices for each night are derived from one or more measurements in the different colours taken within a time interval of $``$ 1 hour. The overall trend was for the spectrum to become hotter (more blue) over the principal five weeks of observation. This suggests, assuming that most of the light comes from an accretion disc, a trend towards increasing disc temperature as the accreting matter diffuses inwards. There was also a brief decrease in B-V & V-I (increase in colour temperature) between 2002 May 1 and May 4 (HJD 396 - 399), followed by a recovery to the overall trend line. This can be seen in Fig. 2 where, for the first few days following optical identification, the I band flux decreased slightly while fluxes in the other bands increased. The light curves in BVRI are approximately triangular in profile, on a linear scale, and similar to, but delayed by, $``$ 13 days relative to the X-ray light curve. During the time interval May 1 to May 11 (HJD 396 - 406) the BVRI fluxes increased by $`75`$ per cent while the X-ray flux decreased by $`50`$ per cent. We can think of no physical process whereby X-ray emission can lead optical emission by 13 days. In SAX J1808.4–3658, the optical decline preceded the X-ray by $`3\pm 1d`$ (Giles et al. 1999). Since we have no information on the optical flux during April we conclude that the apparent similarity between the optical and X-ray light curves is coincidental. The optical decay evident in the centre panel of Fig. 2 has an e-folding time of $`22.2\pm 1.1`$ d. This appears to be qualitatively different to that for SAX J1808.4–3658 and IGR J00291+5934 ### 3.4 Search for binary modulation On the nights of May 2, 3 & 4 (HJD 397, 398 & 399) we monitored the object at I continuously in an attempt to observe an orbital period modulation. These data are plotted in Fig. 4. In mid-May Galloway et al. (2002a) detected a 2614.75(15) s orbital period modulation of the X-ray pulsation frequency. No amplitude modulation of the X-ray flux was detected. Their orbit ephemeris placed the neutron star on the far side of the companion at 2002 May 11.4941(2) UT (HJD 405.9941). The solid sine curves in the lower part of each panel in Fig. 4 represent the modulation (arbitrary amplitude) expected from X-ray heating of the companion star as in SAX J1808.4–3658 (see fig. 1 in Giles et al. 1999 and fig. 2 in Homer et al. 2002). In Fig. 4 we also show the light curves for a nearby ’constant’ comparison star ($`6\mathrm{}`$ to the NNW in Fig. 1, I mag 17.11) which has been shifted by an arbitary amount. This star does not show the modulation evident for XTE J0929–314 in the top panel. All three observing runs terminated close to the telescope elevation limit but the source exhibits significant variability which is not apparent for the plotted comparison star or for others not shown here. The light curve for the last four hours of the night of May 2 (HJD 397) shows clear indications of modulation at or near the orbital period together with slow changes in the average magnitude. We corrected the data commencing from HJD 396.92 ($`22`$h UT) for these slow changes using a second order polynomial and used the Q method (Warner & Robinson, 1972) to search the de-trended data for periodicity in the range 0.01 to 0.05 days. There is a strong single peak corresponding to a period of $`0.030\pm 0.001`$ days. This is consistent with the orbital period of 0.030263 days discovered by Galloway et al. (2002b). In Fig. 5. we plot the corrected light curve folded at the X-ray period and ephemeris where phase zero is defined as the time when the companion is at its greatest distance from the observer i.e. it lies beyond the neutron star. The vertical error bars are those generated by the DoPHOT photometry which appear to be slightly over-estimated. To clarify this we repeated the phase folding process for the nearby comparison star, using the same X-ray ephemeris, so that the real scatter for a ’constant’ source can be seen. These points are plotted in the lower part of Fig. 5 and have a slightly different offset to that used in Fig. 4. The formal DoPHOT average error for these 22 points is 0.021 mag but they have a $`\pm 1\sigma `$ scatter of only 0.017 mag so the XTE J0929–314 error bars in Figs. 4 & 5 are probably $``$ 20 per cent too large. The modulation is approximately sinusoidal with amplitude $`0.09`$ magnitudes peak to peak and maximum at phase $`0.19\pm 0.05`$. Hence the modulation is unlikely to be due to X-ray heating of the companion as this would have maximum light at phase zero. In this respect XTE J0929–314 differs from SAX J1808.4–3658 in which the orbital period optical modulation had a maximum at phase zero (Giles et al., 1999). Nor is it likely that the modulation is due to emission from a hot spot on the disc since this would have a maximum at a phase between $`0.30.5`$. Perhaps both processes contribute. The first PCA X-ray observation (obs. 01-00) occurred during the time we detected an orbital period modulation on the night of 2002 May 2 (HJD 397) and its duration is shown in the top panel of Fig. 4. However, this observation consists of crossed slews to determine the X-ray source position and is thus not suitable for modulation analysis. As noted earlier, no X-ray amplitude modulation was reported in any of the many following RXTE PCA observations (Juett et al. 2003). ### 3.5 Spectral changes There are many occasions within Tables 2 & 3 where we have BVRI values on the same night but some caution is required in combining these into broadband spectra due to the variability detected on several nights. In most instances the different colour averages for each night are derived from one or more measurements taken in a time interval of $``$ 1 hour. It should be clearly noted that the central wavelength and bandwidth for the R & I filters differs between the Cousins and older Johnston systems and this can affect the apparent spectral shape. Here we use the Cousins R & I parameters. In Fig. 6 we plot the broadband BVRI spectra from 8 nights, 2002 May 1, 4, 11, 13, 14, 19 & June 1, 7 (HJD 396, 399, 406, 408, 409, 414, 427 & 432). Also shown is a curve representing a power law approximation to the emission from an optically thick, X-ray heated disc. The distribution is given by the equation $`F_\lambda \lambda ^3e^{A_\lambda /1.086}`$ where $`F_\lambda `$ is the reddened flux at wavelength $`\lambda `$ and $`A_\lambda `$ is the wavelength dependent reddening correction toward the source. The amplitude is arbitrary and the spectrum is reddened assuming interstellar extinction $`A_V=0.42`$. This value is scaled from the estimated value $`A_V=0.68`$ for SAX J1808.4–3658 (Wang et al., 2001) using the integrated column densities, $`N_H1.3\times 10^{20}cm^2`$ for SAX J1808.4–3658 (Gilfanov et al., 1998) and $`N_H7.6\times 10^{20}cm^2`$ for XTE J0929–314 (Juett et al., 2003). A similar value for $`A_V`$ is obtained by using the relationship between $`A_V`$ and $`N_H`$ given by Predehl & Schmitt (1995). We have no information on the flux in B for the first night (HJD 396) but it is clear that the spectrum was heavily reddened on that occasion. On subsequent nights the spectra were generally steeper and had an approximately power law distribution. There was, however, a highly variable red excess above the power law. The excesses are not sensitive to the assumed value of $`A_V`$. On HJD 399 the excess was near zero (similar to the canonical disc power law distribution) and, as noted in section 3.3, the spectrum was anomalously blue. On HJD 409 (and possibly HJD 432 although with less statistical significance) the R band was strongly enhanced consistent with strong Balmer line ($`H\alpha `$) emission. Balmer emission cannot however account for the variable excesses in I. Measurements errors were relatively large for the last two nights in Fig. 6 (HJD 427 and 432) but it is clear that the red excess had disappeared and that the spectra were steeper (bluer) than during earlier observations. ## 4 Discussion Several features distinguish XTE J0929–314 from SAX J1808.4–3658. Firstly, the maximum of the orbital period modulation occurs at phase $`0.19\pm 0.05`$ rather than at phase zero. This points to an origin of X-ray heating in SAX J1808.4–3658 but the situation is more complex in XTE J0929–314. Galloway et al. (2002b) have shown that the companion in XTE J0929–314 is probably a very low mass ($`0.008M_{}`$) helium white dwarf. The companion in SAX J1808.4–3658 is believed to be a $`0.05M_{}`$ brown dwarf (Bildsten & Chakrabarty, 2001). The Roche lobe radii of the companion stars will therefore differ by almost an order of magnitude substantially reducing the X-ray radiation reprocessed by the companion in XTE J0929–314. We have estimated the amplitude $`L_{OA}`$ of optical modulation due to heating of the companion star in XTE J0929–314 using the relation $`L_{OA}=L_{XA}/L_{XB}(d_B/d_A)^2(R_A/R_B)^2L_{OB}`$ where $`L_{OB}`$ is the optical modulation observed for SAX J1808.4–3658 (Giles et al, 1999), $`L_{XA}/L_{XB}`$ is the ratio of the X-ray fluxes (Gilfanov et al., 1998, Juett et al., 2003), $`d_B/d_A`$ is the ratio of distances between the neutron stars and their companions and $`R_A/R_B`$ is the ratio of the Roche Lobe radii for the two systems. The estimated orbital modulation is $`25`$ per cent of that observed suggesting that, if it is generated thermally, much of the emission comes from a hot spot on the disc. This provides an explanation for the observed maximum phase which lies mid-way between that expected from heating of the companion and from hot spot emission. Turning to the spectral characteristics the broadband spectra from SAX J1808.4–3658 are well fitted by smoothly varying functions derived from a thermal disc model (Wang et al., 2001). However, there are signifiant, time dependent, R & I band excesses in our XTE J0929–314 spectra. Nothing like this has been reported for SAX J1808.4–3658 although Wang et al. (2001) reported a strong near IR JHK excess on one occasion and this may well have extended to optical wavelengths. We have insufficient information to explain our observed spectral variability. As noted in section 3.5, variable Balmer emission may be responsible for the R band excesses but cannot contribute to the excesses in I. The remarkable changes in the red excesses between HJD 396 and 399 and between HJD 409 and 414 might perhaps be due to diffusion inwards of cool matter from a brief enhancement of mass transfer onto the disc from the companion star. Alternatively, the red excesses may be due to transient synchrotron emission from matter flowing out of the system via bipolar jets. The synchrotron spectrum is cut off at wavelengths shorter than R. As noted in section 1, many different classes of X-ray binaries emit synchrotron radiation (Fender 2003). Rupen et al. (2002) identified a weak ($`0.35\pm 0.07`$ mJy) 4.86 GHz radio source at the SAX J0929-314 position on 2002 May 3 and 7. Unfortunately there do not appear to have been any IR observations during this outburst. We hypothesise that a rather similar phenomenon may have occurred during the 1998 outburst of the accreting millisecond pulsar SAX J1808.4–3658. Wang et al (2001) noted that the V & I band fluxes measured by the JKT 1-m telescope on 1998 April 18.2 were about 0.2 magnitudes brighter than they were $``$ 0.5 days later when measured by the Mt Canopus 1-m telescope. They assumed that the discrepancy was due to calibration uncertainties in the JKT data. There was a clear IR (JHK) excess measured by the UKIRT telescope at 1998 April 18.6 just 0.4 days after the JKT measurements. Wang et al. (2001) proposed a synchrotron origin for this IR excess. We suggest that the synchrotron excess was also present at the time of the JKT measurements and that it extended into the optical bands. It had disappeared at the time of the Mt Canopus measurements a few hours later. ## 5 Conclusions The optical counterpart of XTE J0929–314 was variable on all timescales down to a few hours during the 2002 May observations. On one occasion lasting $`4`$ hours the I band flux was modulated at the orbital period with amplitude $`0.09\pm 0.01`$ magnitudes. No variability was apparent in the X-ray measurements (Galloway et al, 2002b). The peak of the orbital modulation occurs at a phase of $`0.19\pm 0.05`$ relative to the X-ray ephemeris and appears to rule out X-ray heating of the companion as the source of the modulation unless it is combined with emission from a hot spot on the disc. Broad band BVRI spectra taken on 8 nights have an approximately power law distribution as expected for an optically thick accretion disc but with variable excesses in R & I. Overall these excesses declined and the spectra steepened (became bluer) during the period of the observations. While variable $`H_\alpha `$ emission may be responsible for some of the excess in R, another explanantion is required for the I band enhancements. We suggest they may be due to emission from cool matter in the outer part of the disc following a transient episode of mass transfer from the companion. Alternatively, variable synchrotron emission, cut off at R band wavelengths, contributes to the emission spectrum. There is a clear need for fast follow-up optical, IR and radio observations of millisecond X-ray pulsars. These should include polarimetry and high time resolution optical and near IR photometry to test the synchrotron emission hypothesis. The feasibility of high speed photometry at the pulsar spin frequency might also be investigated. Facilities for immediate data reduction and generation of light curves are essential in order to optimise observing strategies for these highly variable objects. ## 6 Acknowledgements We thank Ron Remillard for timely information on the occurrence of this transient and Duncan Galloway for providing PCA data. This research has made use of data obtained through the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA / Goddard Space Flight Center. We thank Don Melrose, Mark Walker and the referee for helpful comments and gratefully acknowledge financial support for the Mt Canopus Observatory by Mr David Warren. ABG thanks the University of Tasmania Antarctic CRC for the use of computer facilities.
warning/0506/nlin0506046.html	ar5iv	text	# Nonlocal competition and logistic growth: patterns, defects and fronts ## I Introduction Recently, there is a growing interest in the spatial properties of logistic growth with nonlocal interactions Sakai ; Doebeli ; Bolker ; Tokita ; Hoopes ; Anderson ; Sayama ; Fuentes ; Birch ; shnerb ; Garcia . A variety of models have been introduced, including various types of interaction kernels, deterministic and stochastic evolution and growth or death rate that depends on the local population. A common feature found in all these models is the segregation transition, i.e., for small enough diffusion and for certain interaction kernels the homogenous state of the system becomes unstable and the steady state is spatially heterogenous. This feature turns out to be stable against the stochasticity induced by the discrete nature of the reactants, and the total carrying capacity (per unit volume) of the stochastic system depends on the details of the spatial segregation Birch ; Garcia . In previous work shnerb , the general conditions for the integral kernel to allow for spatial segregation have been presented, and the existence of topological defects between ordered domains has been analyzed in detail for a logistic growth on a one dimensional array of patches with nearest-neighbor competition. Here, a comprehensive study of this reaction-diffusion equation is presented: short-range interactions are shown to yield spatial modulation of arbitrary large wavelength and different type of defects, the total population of the system admits nontrivial dependence upon the diffusion rate, and the dynamics of the system is studied, both for global initiation and for local initiation. The appearance of domains with different order parameter and the features of the boundaries between them is considered in detail for various situations. Our starting point is the well-investigated Fisher-KPP equation fisher ; kol , first introduced by Fisher to describe the spread of a favored gene in population: $$\frac{c(x,t)}{t}=D^2c(x,t)+ac(x,t)bc^2(x,t).$$ (1) Clearly, this equation is a straightforward generalization of the logistic growth to spatial domains, and allows for two steady states: an unstable state with $`c(x)=0x`$ and the stable steady state $`c(x)=a/b`$. It was shown that, for any local initiation of the instability (i.e., $`c(x)0`$ on a compact domain) the invasion of the stable phase into the unstable region takes place via a front that moves in a constant velocity $`v_F=2\sqrt{Da}`$. The stability of this solution, the fact that the velocity is determined by the leading edge (”pulled front”) and the corrections to this expression due to stochastic noise associated with the discrete nature of the reactants derrida has been reviewed, recently, by various authors saarloos . The FKPP equation is the simplest equation that describes the transition from unstable to stable steady state on spatial domains, and as such it fits many situations, from the spread of a disease by infection to the advance of a fire or new technology. Accordingly, this model has been widely studied from many points of view and has been generalized in many directions such as modified interaction terms, non linear diffusion and so on. The process considered here, logistic growth with nonlocal competition, is described by the generalized FKPP equation: $$\frac{c(x,t)}{t}=D^2c(x,t)+ac(x,t)c(x,t)_{\mathrm{}}^{\mathrm{}}\gamma (x,y)c(y,t)𝑑y,$$ (2) where $`\gamma (x,y)`$ is the interaction kernel, and the original FKPP process corresponds to the limit $`\gamma (x,y)=\delta (xy)`$. The motivation for the study of this process comes from one of the basic mechanisms in population growth, namely, the competition for common resource. In any autocatalytic system the multiplication of agents depends on various resources (energy, chemicals, water etc.). If there is only limited amount of the resource, its consumption leads to extinction, so generally any crucial resource should be deposited, and its availability dictates the saturation value for the population. As a concrete example let us look at vegetation merron ; lavee : the common resource needed for vegetation is water, and the rain corresponds to deposition of this resource. If the resource dynamics is much faster than that of the agents (shrubs, trees etc.), there is, at any time, a soil moisture profile that reflects the instantaneous vegetation configuration, and there is a depletion of this moisture at the spatial region around a biomass unit. Accordingly, the environmental conditions for a new agent at this region becomes hostile. Following arguments of this type one suggests that competition for common resource induces long range competition among agents via the depletion of the resource in the surroundings of an agent. Another examples may involve the competition for light huisman and cooperation among agents (symbiosis) that may yield ”negative competition” among the reactants. Numerical simulations of the dynamics corresponds to Eq. (2) require space and time discretization. In this work the time evolution of the system is generated via forward Euler integration, where the time step is taken small enough such that further reduction of it do not effect the results. We simulate a system of discrete patches, where the hopping rate is proportional to the diffusion constant. Let us present some a-priori considerations related to this system. There are few basic types of steady state solutions: first, it may happens that the steady state is homogenous: this may be the case if the long range competition is too small, or if the interaction kernel do not allow for the instability to occur shnerb ; Sayama . At some point in the parameter region a bifurcation may occur, and the homogenous states becomes unstable to modulation of wavelength $`2\pi /k`$. Right above the bifurcation one expects, though, to see an inhomogeneous (modulated) steady state. Far from the bifurcation point there are many unstable wavelengthes and some sort of mode competition takes place. In the other limit, i.e., very strong competition, one may expect that ”life” at a single patch forces all the other patches at finite range to be (almost) empty, so the steady state is sort of ”spiky” phase, where many active wavelengthes participate in the formation of localized bumps. As we are looking at a dynamic system with no noise, few stable steady states may exist simultaneously, each admits its own basin of attraction in the space of possible initial conditions. Numerically, however, it turns out that only one important distinction should be made, namely, between local and global initiation: the initiation is ”local” if at $`t=0`$ there is finite support to the colony, while if the system begins with random small biomass that spreads all around it corresponds to global initiation. Within each of these subclasses, the numerics suggests that a generic initial condition will flow into a unique steady state. This paper is organized as follows: in the second section the stable spatial configurations (steady states stable against small fluctuations) are presented: the conditions for an instability of the homogenous solution are reviewed and discussed, and the properties of the final state are identified in different parameter regions, leading to a characteristic ”phase diagram”. In the next section the appearance of defects (separating spatial regions with different order parameter phase) is studied. The fourth section deals with the ”spiky” phase, where many excited modes superimposed to yield a pattern of spikes and the typical defect is a combination of two depletion regions. In the fifth section there is a brief description of phases and defects in two spatial dimensions, and next the effect of the spatial segregation on the global population is considered. In the seventh section the dynamic properties of the model are discussed, including the velocity of the primary and the secondary Fisher fronts and the appearance of topological defect in the invaded region. Some comments and conclusions are presented at the end. ## II Static properties In this section we consider the steady state solutions for equation (2) on spatial domain of coupled, identical patches. The initiation is assumed to be global, i.e., the initial conditions are small, randomly spread, reactant population at each spatial patch. The model considered here allows for nontrivial spatial organization even in the absence of diffusion, due to the long range competition, and global initiation helps to see these features within reasonable simulation times. The differences, if any, between global and local initiation will be considered in the last section. ### II.1 The bifurcation cascade Let us consider the spatially discretized version of (2), i.e., an infinite one dimensional array of identical patches coupled to each other by diffusion and long range competition. The time evolution of the reactant density at the n-th site, $`\stackrel{~}{c}_n`$, is given by: $`{\displaystyle \frac{\stackrel{~}{c}_n(t)}{t}}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{D}}{l_0^2}}[2\stackrel{~}{c}_n(t)+\stackrel{~}{c}_{n+1}(t)+\stackrel{~}{c}_{n1}(t)]+a\stackrel{~}{c}_n(t)`$ (3) $``$ $`b\stackrel{~}{c}_n^2(t)\stackrel{~}{c}_n(t){\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\stackrel{~}{\gamma }_r[\stackrel{~}{c}_{n+r}(t)+\stackrel{~}{c}_{nr}(t)].`$ where $`\stackrel{~}{D}`$ is the diffusion constant and $`a,b,\stackrel{~}{\gamma }`$ are the corresponding reaction coefficients. One may precede to define the dimensionless quantities, $$\tau =at,c=b\stackrel{~}{c}/a,\gamma _r=\stackrel{~}{\gamma }_r/b,D=\frac{\stackrel{~}{D}}{al_0^2}.$$ (4) Note that the new ”diffusion constant” is $`D=W^2/l_0^2`$, where $`W\sqrt{D/a}`$ is the width of the Fisher front, so the dimensionless diffusion is determined by the ratio between the front width and the lattice constant. The continuum limit, though, is the limit where the front width is large in units of lattice spacing. With these definitions Eq. (3) takes its dimensionless form, $`{\displaystyle \frac{c_n}{\tau }}`$ $`=`$ $`D[2c_n+c_{n+1}+c_{n1}]`$ (5) $`+`$ $`c_n\left(1c_n{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\gamma _r[c_{n+r}+c_{nr}]\right),`$ that may be expressed in Fourier space \[with $`A_k_nc_ne^{iknl_0}`$\] as, $$\dot{A}_k=\alpha _kA_k\underset{q}{}\beta _{kq}A_qA_{kq},$$ (6) where $`\alpha _k12D[1cos(kl_0)]`$ (7) $`\beta _k1+2{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\gamma _rcos(rkl_0).`$ (8) Following ns , one observes that $`c_n`$ is positive semi-definite so $`A_0`$ is always ”macroscopic”. Any mode is suppressed by $`A_0`$; accordingly, for small $`\gamma _r`$ one expects only the zero mode to survive. If, on the other hand, $`\gamma _r`$ increases above some threshold, bifurcation may occur with the activization of some other $`k`$ mode(s), and the homogenous solution becomes unstable. The condition for occurrence of such bifurcation is that: $$g(k)=\beta _k+2\beta _0D[1cos(kl_0)]<0$$ (9) fulfilled by some k. This is the situation where patterns appear and translational symmetry breaks. Right above the bifurcation there is only one active $`k`$ mode that dictates the modulation of the system. As $`g(k)`$ decreases further there are many active modes that compete with each other via the nonlinear terms of (6), and the linear stability analysis of the homogenous state may be irrelevant to the final spatial configuration. ### II.2 Nearest neighbors interactions In previous work, the properties of the system have been considered for the extreme case where the competition takes place only between neighboring sites ($`\gamma _r=\gamma `$ for $`r=1`$ AND $`\gamma _r=0`$ if $`r>1`$). For nearest neighbors (n.n) interaction of that type the only stable wave number is $`k=\pi /l_0`$, where $`l_0`$ is the lattice constant, and the bifurcation takes place at $`\gamma <1/2`$. The spatial state at this wave number is $`u_n=A_0+A_\pi cos(n\pi /l_0)`$ and the spatial structure is of the form …ududud… (u=up, large amount of biomass, d=down, small amount). In the absence of diffusion spatial segregation takes the form 101010 , i.e., only the even (odd) sites are populated. Obviously, starting from generic random state different domains are crated with odd or even ”order parameter” and kinks (domain walls) emerge between different domains. As shown in shnerb , the structure of these topological defects, including their size (that diverges at the segregation transition) and their exact form may be calculated analytically. ### II.3 Next nearest neighbors (n.n.n) Quiet surprisingly, the increase of the competition radius by a single site takes us to a completely different regime. While in the case of nearest neighbors interaction the spatial modulation length and the competition length are the same, next nearest neighbors competition (and, accordingly, any interaction of longer range) may yield, upon tuning the parameters, spatial modulation of arbitrary large wavelength. This situation resembles the case of magnetic systems, e.g., an Ising chain: if the exchange interaction is only between nearest neighbors the equilibrium state admits only an up-down modulations, while n.n.n. interaction may yield large solitons, as shown by bak . In that sense the next n.n. case demonstrate the essential features of the long range competition model in a generic way, while at least part of the results may be inferred analytically. The most general form of next nearest neighbors interactions is given by Eq. (5) with: $$\gamma _r=\{\begin{array}{cc}\gamma _1& r=1\\ \gamma _2& r=2\\ 0& \text{else}\end{array}$$ (10) The bifurcation threshold is defined now by the equation: $$g(k)=1+2\gamma _1cos(kl_0)+\gamma _2cos(2kl_0)\beta _k+2\beta _0D[1cos(kl_0)]=0.$$ (11) where $`g`$ has extremum points at $`k_{1,2}=0,\pi `$ and $$k_3=arcos(\frac{\gamma _1+\beta _0D}{4\gamma _2}).$$ (12) If a real wavevector $`k_3`$ exists (i.e., at $`\|\frac{\gamma _1+\beta _0D}{4\gamma _2}\|<1`$) it will be the minimum of $`g(k)`$ while $`k_{1,2}`$ are maxima. For the range of parameters where $`k_3`$ is imaginary the minima may be at $`k=\pi `$ and the modulation is of ”up-down” type, or $`k=0`$, where the homogenous state is stable. The resulting phase diagram, in the $`\gamma _1\gamma _2`$ plane with zero diffusion, is presented in figure (1): In region I the homogenous state is stable, while in region II the bifurcation takes the system to the up-down mode, like the situation for n.n. interaction. In region III, however, $`k_3`$ dominate and modulations of any size may occur. The bifurcation line is given (in the presence of diffusion) by the two branches of the equation: $$\gamma _2=\frac{1+2DD^2+\gamma _1(5D2D^2)\pm \sqrt{1+4D+14D\gamma _12\gamma _1^2+12D\gamma _1^2}}{2D^2+416D}$$ (13) that reduces, at the $`D=0`$ case, to the simple form: $$\gamma _2=\frac{1\pm \sqrt{12\gamma _1^2}}{4}.$$ (14) #### II.3.1 wavelength selection, mode competition and the spiky phase From (12) it seems that the bifurcation wavelength is bounded from above by the interaction length. This, however, is not the actual situation on a discrete lattice: the wavelength inferred from Eq. (12), although bounded, is generically incommensurate with the lattice constant, and the system should choose a commensurate one. It turns out that, if the wavelength is rational (i.e., if $`k_3=2\pi m/n`$, where $`m`$ and $`n`$ admits no common denominator) the spatial modulation repeats itself after $`n`$ lattice sites. A typical example is the steady state obtained numerically for the case $`m=7,n=20`$ where a period-20 modulation appears, as demonstrated in Figure (2). At finite system the maximal $`n`$ allowed is of order of the system size, and only in an infinite system all rational fractions may be activated. Note that in an infinite system any change of the interaction parameters yields different wavelength, a phenomena that resembles the ”devil staircase” situation in spin systems bak . For finite system, though, there is a set of points along the bifurcation line that correspond to the allowed wavelengthes. Numerical simulation indicates that there is a basin of attraction around each of these points, i.e., if the interaction parameters $`\gamma _1`$ and $`\gamma _2`$ yields a prohibited wavelength the system flows into one of the closest allowed modulations. The overall structure is demonstrated in figure (3): close to an isolated point there is a basin of attraction, but further away from the bifurcation line these regions begin to overlap, and the system flows into some mixture of the closest allowed states, depending on its initial conditions. Deep in region III many wavenumbers are involved; the interaction parameters are relatively large, and instead of simple harmonic modulation the system flows, generically, into a spiky steady where the ”living colonies” are separated by the interaction length and are not effected by the competition between patches. Although the numerical examples presented here are for a system with next nearest neighbor competition and without diffusion, it is easy to extract from it the properties of the steady state in general. The effect of diffusion is to increase the size of the stable region so the bifurcation line of Figure (1) moves outward together with the pure and the spiky states. For interactions of longer range the parameter space is of higher dimensionality but all other features are essentially the same. ## III Defects The transition from the homogenous to the modulated state involves spontaneous breakdown of translational symmetry, and upon global initiation one may expect domain walls, or kinks, that separate spatial regions with different order parameter. The presence of these defect and their character is crucial for the understanding of the system response functions, e.g., its behavior under small noise: as there is no preference to one phase of the order parameter the kinks may move freely, while the ”bulk” of the domain is much more stiff. In the following paragraphs the characteristic defects for various phases are presented. ### III.1 Domain walls As mentioned above, the nearest neighbors competition leads, above the bifurcation threshold, to to appearance of an up-down modulation ($`k=\pi `$), and if there is no diffusion the steady state is the 0101010 configuration. Clearly there are two equivalent segregation of this type, namely, filled odd sites and empty even sites and vise versa. Accordingly, in case of global initiation (random ”seeds” are spread all around) one finds domains of the stable patterns with different parity, and domain walls (technically known as kinks or solitons) that separate these regions, as seen in figure (4). The nearest neighbor interaction is simple enough to allow for an analytic solution for the kink, and the numerical results confirm the predictions shnerb . In the presence of diffusion there is a ”smearing” of the above results: the homogenous state is stable for larger $`\gamma `$, and above the segregation threshold the steady state is ”smeared” from $`\mathrm{}01010\mathrm{}`$ to an ”up - down - up -down” form, and the kinks are not of finite size but admit exponentially decaying tails. See shnerb for details. ### III.2 Phase shift Unlike the nearest neighbor case, competition of longer range leads to instabilities with wavelength of more than one site, i.e., $`c_n=A_0+A_kcos(nkl_0)`$ with general $`k`$. This opens the problem of defects between ordered regions. Inspired by the nearest neighbors example one may expect another types of kinks that separate different regions of ordered state. Surprisingly, this is not the case. Instead of getting kinks between different oriented regions of the activated wavenumber, one gets single oriented region with phase shift, namely, the spatial structure is of the form $`c_n=A_0+Bcos(nkl_0+\varphi )`$, where $`\varphi `$ is the phase shift between the actual solution and the predicted modulation $`c_n=A_0+A_kcos(nkl_0)`$ and $`B=A_k/cos\varphi `$. On the unit cycle (Figure ) the meaning of this additional phase is a shift of all point by $`\varphi `$. This shift may reduce the number distinct values in one cycle by one, as indicated in the example of Figure (5): here, instead of six distinct values taken by $`c_n`$ along one wavelength, there are only five. Both numerical simulations of the system dynamics, starting from random initial conditions, and stability analyzes of the possible steady state for arbitrary $`\varphi `$ indicates that, although any $`\varphi `$ corresponds to locally stable solution, the most stable $`\varphi `$ equals to half of the angular distance between two adjusting points on the circle. In Figure (5) the actual phase shifted pattern is shown for $`k=3\pi /5l_0`$, while Figure (6) indicates that the most stable phase corresponds to $`\varphi /\varphi _0=1/2`$. As the Lyapunov exponent of any $`\varphi `$ is negative, small perturbations around any $`\varphi `$ value (in particular, $`\varphi =0`$) decay. Figure (7) shows the corresponding stable mode with $`\varphi =0`$ where the initial conditions are small perturbation around it. Figure (8), on the other hand, shows the final state with generic initial conditions, where the system flows to the most stable pattern with $`\varphi /\varphi _0=1/2`$. ## IV The spiky phase Deep in region III of the phase diagram \[Figure (1)\] many wavevectors are excited, with strong mode competition between them, and the linear analysis picture based on Fourier decomposition becomes ineffective. Better insight into the system comes from a real space analysis: deep into region III the long range competition is strong, and within the effective interaction range new colony can not develop in the presence of a fully grown one. Accordingly, this phase is characterized by fully developed colonies separated by ”dead regions” of constant length that reflects the effective interaction length. In Fourier space, this corresponds to many active modes that build together a periodic structure of ”bumps”. In case of global initiation, of course, defects may appear in the stable steady state as the system flows to different order parameter in different regions. Again, it is better to use the real space picture in order to describe these defects. The situation is close to what observed in the case of random sequential adsorption ads ; lavee : while an ”optimal” filling of the system admits a periodic structure of living patches with periodicity of, say, $`L`$ lattice points, it may happens that the distance between two fully developed sites is between L and 2L, and all the site in between should remain empty due to the long range competition. The emerging spatial configuration is of ordered regions (with coherence size that depends upon the dynamics) separated by ”domain walls”, where the width of these walls is taken from some distribution function between zero and the interaction effective length. ## V Two dimensional system Although all the analysis presented was in one dimension, the basic picture is the same for higher dimensionality. In particular, the bifurcation condition is similar, nearest neighbor interactions yields a ”checkerboard” phase above the bifurcation line, and the spiky phases is also observed. For nearest neighbors interaction kinks between different regions (checkerboard parity) occurs. Because of the two dimensionality of the lattice the kinks might have any arbitrary spatial line, rather then straight line, as shown at figure (9). Those kinks are de-facto one dimensional topological defects, because of the periodic boundary conditions. On the other hand, the domain walls of Figure (10) seems to admit a real 2d features, although their topological character is not clear. ## VI Global properties ### VI.1 Upper critical diffusion Let us turn back to the bifurcation condition, Eq. (9), in different representation: $$\frac{\beta _k}{\beta _0}+2D[1cos(kl_0)]<0$$ (15) where the $`k`$ considered is the one for which $`\beta _k`$ admits a global minimum. Clearly, this $`k_{min}`$ depends only on the form of the interaction kernel and is independent of its strength (if one multiplies all $`\gamma _r`$ by constant factor, the value of $`k_{min}`$ remains the same). Since the negative term in the instability condition $`\frac{\beta _k}{\beta _0}`$ cannot exceed $`(1)`$, the absolute value of the right hand term should be even smaller to allow a periodic modulation of the stable steady state. Assume, now, that the wavelength of the modulation is much larger than the lattice constant (as already required as one approaches the the continuum limit). In that case the approximation $`2D[1cos(kl_0)]Dk_{min}^2l_0^2`$ holds, and since $`D=(W_F/l_0)^2`$, this term is proportional to $`(W_F/\lambda )^2`$, where $`\lambda `$ is the period of the modulation. This implies that, independent of the strength of the long-range competition, bifurcation never takes place if the width of the Fisher front is larger than the period of the modulation. This statement holds up to a numerical factor (between zero and one) which determined by the form of the competition kernel. Simple example that demonstrate these considerations is the case of nearest neighbors interaction. Here $$g(k)=1+2\gamma cos(k)+2(1+2\gamma )D[1cos(kl_0)]$$ (16) and the global minima is $`k=\pi `$. $`g(k_{min})`$ is $$g(\pi )=12\gamma +4(1+2\gamma )D$$ (17) so for any $`\gamma `$ there is an upper critical D $$D_c=\frac{2\gamma 1}{4(1+2\gamma )}$$ (18) above which no bifurcation takes place. This upper critical diffusion constant converges to a global value as $`\gamma \mathrm{}`$ $$D_c^gD_c,_{\gamma >\mathrm{}}=\frac{1}{4}.$$ (19) and no bifurcation takes place when the width of the Fisher front is of order of the modulation length. Intuitively this result may be understood as follows: suppose that the system is in its 010101 state, and suppose that the dynamics is discrete in time. If $`D=1/4`$ it implies that each filled site contribute $`1/4`$ to any of its neighbors, and then the system is frozen in its homogenous state with amplitude $`1/2`$ at each site. Generalizing this intuition to periodic modulation of arbitrary wavelength yields the same result, where the Fisher front width stands as a definition of an ”effective site”. ### VI.2 Spatial segregation and total population Given a system with long range competition, one may ask how the total population (integrated over all the spatial domain) or the average population density, depend on the phase of the system. Naively one expects the total population to grow with the diffusion constant, as faster spatial wandering helps an individual reactant to escape from the depleted region of an already existing colony. This, however, is not the case, as pointed out by Birch and Garcia : the size of the total population depends on the efficiency of segregation: strong segregation implies higher population (on average, since there are empty regions and living patches). Thus, decrease of diffusion implies higher total population density. Clearly, the total population is given by the amplitude of the zero mode in Fourier decomposition of the population, (See Eq. 6). As long as the system is in its homogenous phase this quantity is diffusion independent and the total population depends only on the strength of the interaction, $`A_0=1/\beta _0`$. Right above the bifurcation, when only one excited mode ($`k`$) exists, the total population is proportional to $`A_0=\alpha _k/(\beta _0+\beta _k)`$, and since $`\alpha _k`$ increases as $`D`$ decreases, so is the total population. In the case of one dimensional lattice with nearest neighbors interaction, for example, the dependence of the total occupancy of the sample on the diffusion constant may be calculated explicitly, since there is only one excited mode $`k=\pi /l_0`$. Here even far from the bifurcation point the amplitude of the zero mode is given by $`A_0=\alpha _k/(\beta _0+\beta _k)`$. The total sum versus diffusion is, accordingly, $$A_0=\{\begin{array}{cc}(14D)/2& D<D_c\\ 1/(1+2\gamma )& D>D_c.\end{array}$$ (20) Figure (11) shows the total sum versus diffusion for few situations. The numerical results indicate that the decay of average population is approximately linear. Note that, for the ”top hat” competition presented here, there seems to be a discontinuity at $`D_c`$ in two dimensions, while in $`1d`$ the total population is continuous at the transition. ## VII Local initiation: the dynamics of invasion and segregation Along this paper, an analysis of the stable steady states of the logistic growth with long range competition was presented. As few stable steady solutions may exist simultaneity for the same set of parameters, the generic situation was identified numerically using global initiation, i.e., small random population at each site. In this section, the dynamics of growth is analyzed, where the initial conditions are a colony with compact support. For local logistic growth this problem was considered years ago by Fisher fisher and Kolomogorov kol . The invasion of the stable solution into the unstable one takes place via a front (the Fisher front) that propagates in constant velocity. This problem was considered by many authors in different contexts and was generalized to other cases of invasion into unstable state, see comprehensive review by van Saarloos saarloos . As emphasized above, the system considered here may admit \[in region II and III of figure (1)\] two instabilities: the empty state is unstable against the homogenous one, while the homogenous solution breaks and yields a spatial modulation. Accordingly, if the system is initiated locally from a small colony of compact support one expects that two fronts propagate into the empty region: first the front associated with the homogenous state, and then the modulation (secondary instability) front hohenberg . These two fronts travel in different velocities. Generally, it is known that the Fisher velocity is determined by the leading edge (”pulled” fronts) and is related to the Lyapunov exponent that characterizes the relevant instability. Accordingly, the dynamics of our system is determined by two velocities: $`v_p`$, the velocity of the primary front (that interpolates between the empty and the homogenous state) and the modulation velocity $`v_s`$. While $`v_p`$ is $`\gamma `$ independent, the secondary front velocity $`v_s`$ depends on the characteristics of the long range competition. By tuning of $`\gamma `$, though, one may change the relative velocity between the primary and the secondary front. Both velocities may be calculated analytically using a saddle point method and taking into account the discreteness of the lattice points, as discussed in Appendix A. Generically, there are two possible scenarios for the takeover of an empty region by spatially modulated steady state: in the first case $`v_p>v_s`$ (see Figure 12) and the homogenous region between the primary and the secondary front grows linearly in time. This situation is very sensitive, as small perturbations (induced by the leading front) lead to spontaneous bifurcation of the homogenous region, a process that yields many structural defects (e.g., kinks) along the chain. In the second case the situation is different: if $`v_p<v_s`$ there is no homogenous region, and only one front exists. Its velocity is determined, of course, by the primary front velocity, but its shape is different (see Figure 13). In that case the sensitive homogenous region never exists, and the pattern formation process is robust, with no defects associated with the front kinetics. ## VIII Conclusions and remarks This paper attempt to present the various phases associated with the steady states of the logistic process on spatial domains with nonlocal competition. The main feature is, of course, the segregation transition that happens, as was shown, where the width of the Fisher front (associated with the homogenous solution) becomes shorter than the instability wavelength. Right above the bifurcation one finds a pattern dominated by a single wavelength, while far away from the bifurcation line the stable steady state becomes spiky. Each phase is associated with its own defects: phase shift close to the bifurcation, empty regions in the spiky phase, and domain walls (kinks) for the up-down phase of the nearest-neighbors interaction. It turns out that the segregation transition increase the overall carrying capacity per unit volume. In $`1d`$ the population is continues at the transition while in two dimensions discontinuity might occur. Upon local initiation the system dynamics is governed by the relations between the velocities of the primary (empty to homogenous) and the secondary (homogenous to modulated) fronts. The numerics suggests that, while global initiation may yield ”disordered” structure with many defects per unit length, local initiation with the same parameters yields ordered structure unless the secondary front velocity is smaller that the primary one. While in this work only rate equations of reaction-diffusion type has been considered, in recent numerical works of Birch and Young Birch and Garcia et. al. Garcia the stochastic motion of the individual reactants is taken into account. These stochastic models add two ingredients to the description presented here. First, the introduction of individual reactants (”Brownian bugs” young2 ) implies a threshold on the reactant concentration on single patch. Secund, there is a multiplicative noise associated with the stochastic motion of individual reactants. As shown in this work, many of the features associated with long range competition are independent of the discrete nature of individual reactants. ## IX acknowledgements The authors thank Prof. David Kessler for many helpful discussions. This work was supported by the Israeli Science Foundation, grant no. 281/03, and by Yeshaya Horowitz Fellowship. ## X Appendix A In this appendix the analytic expression for the secondary front velocity on a discrete lattice is obtained, via the saddle point argument (see levin ). For the sake of simplicity, only the case of nearest neighbors interaction is considered. In order to preform the same calculations for competition beyond the n.n. limit, one should first find numerically the steady state modulation and then follow the same procedure. The evolution of a population is given by: $$\frac{c_n}{t}=D[2c_n+c_{n+1}+c_{n1}]+c_nc_n^2+c_n\gamma (c_{n+1}+c_{n1}).$$ (21) Denoting by $`\delta _n`$ the deviations from the homogenous solution, $`c_n=A_0+\delta _n`$, Eq. (21) is linearized to yield: $$\frac{\delta _n}{t}=\alpha \delta _n+\beta (\delta _{n+1}+\delta _{n1})$$ (22) Where $`\alpha =a2bA_0+2A_0\gamma 2D/l_0^2`$ and $`\beta =D/l_0^2A_0\gamma `$. Assuming modulated solution of the form, $$\delta _n=\{\begin{array}{cc}Ae^{ikl_0n+\mathrm{\Gamma }(k)t}& nodd\\ Be^{ikl_0n+\mathrm{\Gamma }(k)t}& neven\end{array}$$ (23) and plugging (23) into (22) one gets $$\mathrm{\Gamma }(k)\left[\begin{array}{c}A\\ B\end{array}\right]=\left[\begin{array}{cc}\alpha & \beta \mathrm{cos}(kl_o)\\ \beta \mathrm{cos}(kl_o)& \alpha \end{array}\right]\left[\begin{array}{c}A\\ B\end{array}\right].$$ (24) The dispersion relations is given by: $$\mathrm{\Gamma }(k)=\alpha +\beta \mathrm{cos}(kl_0).$$ (25) where the plus sign is chosen for the unstable modes. The solutions are of the form $$\left[\begin{array}{c}c_n\\ c_{n+1}\end{array}\right]=\left[\begin{array}{c}A\\ B\end{array}\right]e^{ikx+\mathrm{\Gamma }(k)t}.$$ (26) If a solution represents a travelling front with velocity $`v`$ it is useful to define the coordinate system in the moving frame, $`\zeta =xvt`$, to get $$\left[\begin{array}{c}c_n\\ c_{n+1}\end{array}\right]=\left[\begin{array}{c}A\\ B\end{array}\right]e^{ik\zeta +ikvt+\mathrm{\Gamma }(k)t}.$$ (27) Using the saddle point method saarloos the two equations that determine the velocity are $$fivk+\alpha +2\beta cosh(kl_0)=0$$ (28) and $$\frac{f}{k}=iv+2\beta _0sinh(kl_0)=0.$$ (29) In case of finite time steps one should replace $`ivk`$ by $`(e^{ikvdt}1)/dt`$ to get the appropriate corrections. Figure (14) shows the perfect fit between the solution of (28) and (29) and the numerical solution. ## XI Acknowledgements We thank Prof. David Kessler for many helpful discussions and comments. This work was supported by the Israel Science Foundation (grant no. 281/03) and by the Yeshaya Horowitz Fellowship.
warning/0506/hep-ph0506313.html	ar5iv	text	# Smoking-gun signatures of little Higgs models ## 1 Introduction Elucidating the mechanism of electroweak symmetry breaking (EWSB) is the central goal of particle physics today. A full understanding of EWSB will include a solution to the hierarchy or naturalness problem – that is, why the weak scale is so much lower than the Planck scale. Whatever is responsible for EWSB and its hierarchy, it must manifest experimentally at or below the TeV energy scale. A wide variety of models have been introduced over the past three decades to address EWSB and the hierarchy problem: supersymmetry, extra dimensions, strong dynamics leading to a composite Higgs boson, and the recent “little Higgs” models in which the Higgs is a pseudo-Goldstone boson. In this paper we consider this last possibility. In the little Higgs models, the Standard Model (SM) Higgs doublet appears as a pseudo-Goldstone boson of an approximate global symmetry that is spontaneously broken at the TeV scale. The low energy degrees of freedom are described by nonlinear sigma models, with a cutoff at an energy scale one loop factor above the spontaneous symmetry breaking scale. Thus the little Higgs models require an ultraviolet (UV) completion at roughly the 10 TeV scale. The explicit breaking of the global symmetry, by gauge, Yukawa and scalar interactions, gives the Higgs a mass and non-derivative interactions, as required of the SM Higgs doublet. The little Higgs models are constructed in such a way that no single interaction breaks all of the symmetry forbidding a mass term for the SM Higgs doublet. This collective symmetry breaking guarantees the cancellation of the one-loop quadratically divergent radiative corrections to the Higgs boson mass. Quadratic sensitivity of the Higgs mass to the cutoff scale then arises only at the two-loop level, so that a Higgs mass at the 100 GeV scale, two loop factors below the 10 TeV cutoff, is natural. Little Higgs models can thus stabilize the “little hierarchy” between the electroweak scale and the 10 TeV scale at which strongly-coupled new physics is allowed by electroweak precision constraints. Little Higgs models contain new gauge bosons, a heavy top-like quark, and new scalars, which cancel the quadratically divergent one-loop contributions to the Higgs boson mass from the SM gauge bosons, top quark, and Higgs self-interaction, respectively. Thus the “smoking gun” feature of the little Higgs mechanism is the existence of these new gauge bosons, heavy top-like quark, and new scalars, with the appropriate couplings to the Higgs boson to cancel the one-loop quadratic divergence. Since the little Higgs idea was introduced , many explicit models have been constructed. Since the little Higgs idea could be implemented in a number of ways, it is crucial to pick out the experimental signatures that identify the little Higgs mechanism in addition to those that identify the particular little Higgs model. Detailed phenomenological and experimental studies of little Higgs physics at the CERN Large Hadron Collider (LHC) have so far been carried out only within the “Littlest Higgs” model .<sup>1</sup><sup>1</sup>1The LHC phenomenology of the Littlest Higgs model with $`T`$-parity was studied in Ref. ; models with $`T`$-parity will be briefly discussed in Sec. 2. Fortunately, this effort need not be repeated for each of the many little Higgs models, because the models can be grouped into two classes that share many phenomenological features, including the crucial “smoking gun” signatures that identify the little Higgs mechanism. In this paper we categorize the little Higgs models into two classes based on the structure of the extended electroweak gauge group: models in which the SM SU(2)<sub>L</sub> gauge group arises from the diagonal breaking of two or more gauge groups, called “product group” models , and models in which the SM SU(2)<sub>L</sub> gauge group arises from the breaking of a single larger gauge group down to an SU(2) subgroup, called “simple group” models . (This categorization and nomenclature was introduced in Ref. .) These two classes of models also exhibit an important difference in the implementation of the little Higgs mechanism in the fermion sector. As representatives of the two classes, we study the Littlest Higgs model and the SU(3) simple group model , respectively. We find that by examining the properties of the new heavy fermion(s), one can distinguish the structure of the top quark mass generation mechanism and test the little Higgs mechanism in the top sector. Furthermore, by measuring the couplings of the new TeV-scale gauge bosons to the Higgs, SM gauge bosons, and fermions, one can determine the gauge structure of the extended theory and test the little Higgs mechanism in the gauge sector. To emphasize the “smoking gun” nature of the signals, we also compare our results with other models that give rise to similar signatures. For the heavy top partner, we compare the little Higgs signatures with the signatures of a fourth generation top-prime and of the top quark see-saw model. For the TeV-scale gauge bosons, we compare with the $`Z^{}`$ signatures in $`E_6`$, left-right symmetric, and sequential $`Z^{}`$ models. In each case, we point out the features of the little Higgs model that distinguish it from competing interpretations. The rest of this paper is organized as follows. In the next section we describe the basic features of the two representative models. Specific little Higgs models that fall into each of the two classes are surveyed in Appendix A. In Sec. 3, we discuss the top quark mass generation and the quadratic divergence cancellation mechanism in the two classes of models, describe the resulting differences in phenomenology, and show how to test the little Higgs mechanism in the top sector. We also comment on the phenomenological differences between little Higgs models and other models with extended top sectors. In Sec. 4, we discuss the gauge sectors in the two classes of models and identify features common to the models in each class. We discuss techniques for determining the structure of the extended gauge sector and for testing the little Higgs mechanism in the gauge sector. In Sec. 5 we collect some additional features of the phenomenology of the SU(3) simple group model. We conclude in Sec. 6. Technical details of the SU(3) simple group model are given in Appendix B. ## 2 Two classes of little Higgs models If the little Higgs mechanism is realized in nature, it will be of ultimate importance to verify it at the LHC, by discovering the predicted new particles and determining their specific couplings to the SM fields that guarantee the cancellation of the Higgs mass quadratic divergence. The most important characteristics of implementations of the little Higgs idea are ($`i`$) the structure of the extended gauge symmetry and its breaking pattern, and ($`ii`$) the treatment of the new heavy fermion sector necessary to cancel the Higgs mass quadratic divergence coming from the top quark. As we will see, the distinctive features of both the gauge and top sectors of little Higgs models separate naturally into the product group and simple group classes. The majority of little Higgs models are product group models. In addition to the Littlest Higgs, these include the theory space models (the Big Moose and the Minimal Moose ), the SU(6)/Sp(6) model of Ref. , and two extensions of the Littlest Higgs with built-in custodial SU(2) symmetry . The product group models have the following generic features. First, the models all contain a set of SU(2) gauge bosons at the TeV scale, obtained from the diagonal breaking of two or more gauge groups down to SU(2)<sub>L</sub>, and thus contain free parameters in the gauge sector from the independent gauge couplings. Second, since the collective symmetry breaking in the gauge sector is achieved by multiple gauged subgroups of the global symmetry, models can be built in which the SM Higgs doublet is embedded within a single non-linear sigma model field; many product group models make this simple choice. Third, the fermion sector of this class of models can usually be chosen to be very simple, involving only a single new vector-like quark. The simplest incarnation of the product group class is the so-called Littlest Higgs model , which we briefly review here. It features a \[SU(2)$`\times `$U(1)\]<sup>2</sup> gauge symmetry<sup>2</sup><sup>2</sup>2Strictly speaking, it is not necessary to gauge two factors of U(1) in order to stablize the little hierarchy, because the hypercharge gauge coupling is rather small and does not contribute significantly to the Higgs mass quadratic divergence below a scale of several TeV. Thus, there is an alternate version of the Littlest Higgs model in which only SU(2)$`{}_{}{}^{2}\times `$U(1)<sub>Y</sub> is gauged. embedded in an SU(5) global symmetry. The gauge symmetry is broken by a single vacuum condensate $`f`$ TeV down to the SM SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> gauge symmetry. The SM Higgs doublet is contained in the resulting Goldstone bosons, whose interactions are parameterized by a nonlinear sigma model. The gauge and Yukawa couplings radiatively generate a Higgs potential and trigger EWSB. The new heavy quark sector in the Littlest Higgs model consists of a pair of vectorlike SU(2)-singlet quarks that couple to the top sector. The Lagrangian is $$_Y=\frac{i}{2}\lambda _1fϵ_{ijk}ϵ_{xy}\chi _i\mathrm{\Sigma }_{jx}\mathrm{\Sigma }_{ky}u_3^c+\lambda _2f\stackrel{~}{t}\stackrel{~}{t}^c+\mathrm{h}.\mathrm{c}.,$$ (1) where $`\chi _i=(b_3,t_3,i\stackrel{~}{t})`$ and the factors of $`i`$ in Eq. (1) and $`\chi _i`$ are inserted to make the masses and mixing angles real. The summation indices are $`i,j,k=1,2,3`$ and $`x,y=4,5`$, and $`ϵ_{ijk}`$, $`ϵ_{xy}`$ are antisymmetric tensors. The vacuum expectation value (vev) $`\mathrm{\Sigma }\mathrm{\Sigma }_0`$ marries $`\stackrel{~}{t}`$ to a linear combination of $`u_3^c`$ and $`\stackrel{~}{t}^c`$, giving it a mass of order $`f`$ TeV. The resulting new charge 2/3 quark $`T`$ is an isospin singlet up to its small mixing with the SM top quark (generated after EWSB). The orthogonal linear combination of $`u_3^c`$ and $`\stackrel{~}{t}^c`$ becomes the right-handed top quark and marries $`t_3`$. The scalar interactions of the up-type quarks of the first two generations can be chosen to take the same form as Eq. (1), except that there is no need for an extra $`\stackrel{~}{t}`$, $`\stackrel{~}{t}^c`$ since the contribution to the Higgs mass quadratic divergence from quarks other than top is numerically insignificant below the nonlinear sigma model cutoff $`\mathrm{\Lambda }4\pi f10`$ TeV. In contrast, the simple group models share two features that distinguish them from the product group models. First, the simple group models all contain an SU($`N`$)$`\times `$U(1) gauge symmetry that is broken down to SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub>, yielding a set of TeV-scale gauge bosons. The two gauge couplings of the SU($`N`$)$`\times `$U(1) are fixed in terms of the two SM SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> gauge couplings, leaving no free parameters in the gauge sector once the symmetry-breaking scale is fixed. This gauge structure also forbids mixing between the SM $`W^\pm `$ bosons and the TeV-scale gauge bosons, again in contrast to the product group models. Second, in order to implement the collective symmetry breaking, simple-group models require at least two sigma-model multiplets. The SM Higgs doublet is embedded as a linear combination of the Goldstone bosons from these multiplets. This introduces at least one additional model parameter, which can be chosen as the ratio of the vevs of the sigma-model multiplets. Moreover, due to the enlarged SU($`N`$) gauge symmetry, all SM fermion representations have to be extended to transform as fundamental (or antifundamental) representations of SU($`N`$), giving rise to additional heavy fermions in all three generations. The existence of multiple sigma-model multiplets generically results in a more complicated structure for the fermion couplings to scalars. On the other hand, the existence of heavy fermion states in all three generations as required by the enlarged gauge symmetry provides extra experimental observables that in principle allow one to disentangle this more complicated structure. The simplest incarnation of the simple group class is the SU(3) simple group model . We briefly review its construction here; additional details are presented in Appendix B. The electroweak gauge structure is SU(3)$`\times `$U(1)<sub>X</sub>. There are two sigma-model fields, $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$, transforming as $`\mathrm{𝟑}`$s under SU(3). Vacuum condensates $`\mathrm{\Phi }_{1,2}=(0,0,f_{1,2})^T`$ break SU(3)$`\times `$U(1)<sub>X</sub> down to the SM SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub>. The TeV-scale gauge sector consists of an SU(2)<sub>L</sub> doublet $`(Y^0,X^{})`$ of gauge bosons corresponding to the broken off-diagonal generators of SU(3), and a $`Z^{}`$ gauge boson corresponding to the broken linear combination of the $`T^8`$ generator of SU(3) and the U(1)<sub>X</sub>. The model also contains a singlet pseudoscalar $`\eta `$. The top quark mass is generated by the Lagrangian $$_Y=i\lambda _1^tu_1^c\mathrm{\Phi }_1^{}Q_3+i\lambda _2^tu_2^c\mathrm{\Phi }_2^{}Q_3,$$ (2) where $`Q_3^T=(t,b,iT)`$ and the factors of $`i`$ in Eq. (2) and $`Q_3`$ are again inserted to make the masses and mixing angles real. The $`\mathrm{\Phi }`$ vevs marry $`T`$ to a linear combination of $`u_1^c`$ and $`u_2^c`$, giving it a mass of order $`f`$ TeV. The new charge 2/3 quark $`T`$ is a singlet under SU(2)<sub>L</sub> up to its small mixing with the SM top quark (generated after EWSB). The orthogonal linear combination of $`u_1^c`$ and $`u_2^c`$ becomes the right-handed top quark. For the rest of the quarks, the scalar interactions depend on the choice of their embedding into SU(3). The most straightforward choice is to embed all three generations in a universal way, $`Q_m^T=(u,d,iU)_m`$, so that each quark generation contains a new heavy charge 2/3 quark. This embedding leaves the SU(3) and U(1)<sub>X</sub> gauge groups anomalous; the anomalies can be canceled by adding new spectator fermions at the cutoff scale $`\mathrm{\Lambda }4\pi f`$. An alternate, anomaly-free embedding puts the quarks of the first two generations into antifundamentals of SU(3), $`Q_m^T=(d,u,iD)_m`$, with $`m=1,2`$, so that the first two quark generations each contain a new heavy charge $`1/3`$ quark. Interestingly, an anomaly-free embedding of the SM fermions into SU(3)$`{}_{c}{}^{}\times `$SU(3)$`\times `$U(1)<sub>X</sub> is only possible if the number of generations is a multiple of three .<sup>3</sup><sup>3</sup>3This rule can be violated in models containing fermion generations with non-SM quantum numbers, e.g., mirror families . Electroweak precision observables provide strong constraints on any extensions of the SM. The constraints on the little Higgs models have been studied extensively . Of course, any phenomenological study of a particular model must take these constraints into account. However, in this paper we study the generic phenomenology of classes of little Higgs models, using specific models only as prototypes. We focus on features of the phenomenology that are expected to persist in all models within a given class, in spite of variations in the model that can give rise to very different constraints from electroweak precision observables. For exmaple, variations of the model that improve the electroweak fit will not in general change the generic features of the new heavy top-partner phenomenology. Thus, in order to maintain applicability to a wide range of models in each class, we will not limit our presentation of results to the parameter space allowed by electroweak precision fits in the specific models under consideration. For completeness, we now briefly summarize the results of electroweak precision fits in the models under consideration. The most up-to-date studies are Refs. , which include LEP-2 data above the $`Z`$ pole. In most little Higgs models, particularly the product group models, the electroweak data mostly set lower bounds on the masses of the heavy vector bosons due to their contributions to four-Fermi operators and their mixing with the $`W`$ and $`Z`$ bosons. On the other hand, the most important contributions to the Higgs mass quadratic divergence cancellation come from the top quark partner $`T`$, which should be as light as possible to minimize the fine-tuning. These competing desires dictate the favored parameter regions of the little Higgs models. * Littlest Higgs model: The Littlest Higgs model with \[SU(2)$`\times `$U(1)\]<sup>2</sup> gauged contains a new U(1) boson, $`A_H`$, which is relatively light and tends to give rise to large corrections to electroweak precision observables. Assigning the fermions to transform under SU(2)<sub>1</sub> and U(1)<sub>1</sub> only, Ref. finds a stringent constraint $`f5`$ TeV. However, allowing the fermions to transform under both U(1) groups (as required in order to write down gauge invariant Yukawa couplings in a straightforward way) tends to reduce this constraint; Refs. , which do not include LEP-2 data in their fit, found the constraint on $`f`$ reduced from 4 TeV to about 1 TeV; similarly, Ref. found the constraint reduced from 5 TeV to about 2–3 TeV. Gauging only SU(2)$`{}_{}{}^{2}\times `$U(1)<sub>Y</sub>, Ref. found that $`f>\mathrm{max}(6.5c^2,3.7c)`$ TeV \[$`c`$ is defined below Eq. (34)\]. Thus, for example, $`f>1`$ TeV for $`c1/3`$; this yields a lower bound on the heavy gauge boson mass of $`M_{W_H}=M_{Z_H}2`$ TeV. The mass of the $`T`$ quark is constrained to be $`M_T\sqrt{2}f`$, or in this most favorable case $`M_T1.4`$ TeV. * SU(3) simple group model: Reference expands on the analysis of Ref. for this model by including the effect of the TeV-scale fermions in the universal fermion embedding. For our choice of parameterization, the constraint on $`f\sqrt{f_1^2+f_2^2}`$ is relaxed by going to $`t_\beta f_2/f_1>1`$ . For $`t_\beta =3`$, $`f3.9`$ TeV , corresponding to $`M_Z^{}2.2`$ TeV. The mass of the $`T`$ quark in this model is bounded by $`M_Tf\mathrm{sin}2\beta `$; this constraint then translates into $`M_T2.3`$ TeV. Reference found that the anomaly-free fermion embedding is somewhat favored over the universal embedding by electroweak precision constraints. Finally, we mention briefly a different approach to alleviating the electroweak precision constraints on little Higgs models. Because the little Higgs mechanism for canceling the quadratically divergent radiative corrections to the Higgs mass operates at one-loop, it is possible to impose an additional symmetry, dubbed $`T`$-parity , under which the new gauge bosons and scalars are odd. This eliminates tree-level contributions of the new particles to electroweak precision observables, thereby essentially eliminating the electroweak precision constraints<sup>4</sup><sup>4</sup>4Although $`T`$-parity suppresses the contributions of heavy gauge bosons and heavy top partners to electroweak oblique parameters, there is a contribution to four fermion operators through a box diagram involving mirror fermions and Goldstone bosons that is not suppressed by the same mechanism and does not decouple as the mirror fermions become heavy. The mirror fermions must be kept light (i.e., be introduced into the low energy spectrum) in order to suppress the relevant couplings .. It also changes the collider phenomenology drastically, by eliminating signals from single production of the new particles that are odd under $`T`$-parity: in particular, the heavy gauge bosons can only be produced in pairs, eliminating the distinctive Drell-Yan signal. The heavy top-partners remain even under $`T`$-parity, however, so that their signals are robust. It was shown in Ref. how to add $`T`$-parity to any product group little Higgs model. Ref. also concluded that in simple group models, one cannot find a consistent definition of $`T`$-parity under which all heavy gauge bosons are odd. ## 3 The heavy quark sector The SM top quark gives rise to the largest quadratically divergent correction to the Higgs mass. A characteristic feature of all little Higgs models is the existence of new TeV-scale quark state(s) with specific couplings to the Higgs so that the loops involving the TeV-scale quark(s) cancel the quadratic divergence from the SM top quark loop. Therefore, we begin with a study of the extended top sector of little Higgs models. ### 3.1 Top sector masses and parameters The masses of the top quark $`t`$ and its heavy partner $`T`$ are given in terms of the model parameters by $`m_t=\lambda _tv=\{\begin{array}{cc}{\displaystyle \frac{\lambda _1\lambda _2}{\sqrt{\lambda _1^2+\lambda _2^2}}}v\hfill & \mathrm{in}\mathrm{the}\mathrm{Littlest}\mathrm{Higgs}\mathrm{model},\hfill \\ {\displaystyle \frac{\lambda _1\lambda _2}{\sqrt{2}\sqrt{\lambda _1^2c_\beta ^2+\lambda _2^2s_\beta ^2}}}v\hfill & \mathrm{in}\mathrm{the}\mathrm{SU}(3)\mathrm{simple}\mathrm{group}\mathrm{model};\hfill \end{array}`$ (5) $`M_T=\{\begin{array}{cc}\sqrt{\lambda _1^2+\lambda _2^2}f=(x_\lambda +x_\lambda ^1){\displaystyle \frac{m_t}{v}}f\hfill & \mathrm{in}\mathrm{the}\mathrm{Littlest}\mathrm{Higgs}\mathrm{model},\hfill \\ \sqrt{\lambda _1^2c_\beta ^2+\lambda _2^2s_\beta ^2}f=\sqrt{2}{\displaystyle \frac{t_\beta ^2+x_\lambda ^2}{(1+t_\beta ^2)x_\lambda }}{\displaystyle \frac{m_t}{v}}f\hfill & \mathrm{in}\mathrm{the}\mathrm{SU}(3)\mathrm{simple}\mathrm{group}\mathrm{model}.\hfill \end{array}`$ (8) Fixing the top quark mass $`m_t`$ leaves two free parameters in the Littlest Higgs model, which can be chosen to be $`f`$ and $`x_\lambda \lambda _1/\lambda _2`$. We see that the SU(3) simple group model contains one additional parameter, $`t_\beta \mathrm{tan}\beta =f_2/f_1`$. In the SU(3) simple group model, we define $`f\sqrt{f_1^2+f_2^2}`$. To reduce fine-tuning in the Higgs mass, the top-partner $`T`$ should be as light as possible. The lower bound on $`M_T`$ is obtained for certain parameter choices: $`M_T\{\begin{array}{cc}2{\displaystyle \frac{m_t}{v}}f\sqrt{2}f\hfill & \mathrm{for}x_\lambda =1\mathrm{in}\mathrm{the}\mathrm{Littlest}\mathrm{Higgs}\mathrm{model},\hfill \\ 2\sqrt{2}s_\beta c_\beta {\displaystyle \frac{m_t}{v}}ff\mathrm{sin}2\beta \hfill & \mathrm{for}x_\lambda =t_\beta \mathrm{in}\mathrm{the}\mathrm{SU}(3)\mathrm{simple}\mathrm{group}\mathrm{model},\hfill \end{array}`$ (11) where in the last step we used $`m_t/v1/\sqrt{2}`$. The $`T`$ mass can be lowered in the SU(3) model for fixed $`f`$ by choosing $`t_\beta 1`$, thereby introducing a mild hierarchy between $`f_1`$ and $`f_2`$. With our parameter definitions, the choice $`t_\beta >1`$ reduces the mixing between the light SM fermions and their TeV-scale partners, thereby reducing constraints from $`W`$ coupling universality. ### 3.2 Heavy $`T`$ couplings to Higgs and gauge bosons The couplings of the Higgs doublet to the $`t`$ and $`T`$ mass eigenstates can be written in terms of an effective Lagrangian, $$_Y\lambda _tHt^ct+\lambda _THT^ct+\frac{\lambda _T^{}}{2M_T}HHT^cT+\mathrm{h}.\mathrm{c}.,$$ (12) where the four-point coupling arises from the expansion of the nonlinear sigma model field. This effective Lagrangian leads to three diagrams contributing to the Higgs mass corrections at one-loop level, shown in Fig. 1: (a) the SM top quark diagram, which depends on the well-known SM top Yukawa coupling $`\lambda _t`$; (b) the diagram involving a top quark and a top-partner $`T`$, which depends on the $`HTt`$ coupling $`\lambda _T`$; and (c) the diagram involving a $`T`$ loop coupled to the Higgs doublet via the dimension-five $`HHTT`$ coupling. The couplings in the three diagrams of Fig. 1 must satisfy the following relation in order for the quadratic divergences to cancel: $$\lambda _T^{}=\lambda _t^2+\lambda _T^2.$$ (13) This equation embodies the cancellation of the Higgs mass quadratic divergence in any little Higgs theory. It is of course satisfied by the couplings in both the Littlest Higgs and the SU(3) simple group models, as can be seen by plugging in the explicit couplings given in Table 1. Note that in the SU(3) simple group model, $`\lambda _T`$ vanishes when $`x_\lambda =1`$. If the little Higgs mechanism is realized in nature, it will be of fundamental importance to establish the relation in Eq. (13) experimentally. After EWSB, the coupling $`\lambda _T`$ induces a small mixing of electroweak doublet into $`T`$, $$T=T_0\delta _Tt_0,\delta _T=\lambda _T\frac{v}{M_T},$$ (14) where $`T_0,t_0`$ stand for the electroweak eigenstates before the mass diagonalization at the order of $`v/f`$. This mixing gives rise to the couplings of $`T`$ to the SM states $`bW`$ and $`tZ`$ with the same form as the corresponding SM couplings of the top quark except suppressed by the mixing factor $`\delta _T`$. The Feynman rules are given in Table 1. ### 3.3 Additional heavy quark couplings in the SU(3) simple group model Expanding the SU(2)<sub>L</sub> gauge symmetry to SU(3) forces the introduction of a heavy partner associated with each SU(2)<sub>L</sub> fermion doublet of the SM. The first two generations of quarks are therefore enlarged to contain two new TeV-scale quarks $`Q_{1,2}`$. We consider both the universal and the anomaly-free fermion embeddings, as discussed in more detail in Sec. B.2. The universal embedding gives rise to two charge 2/3 quarks, $`U`$ and $`C`$, while the anomaly-free embedding gives rise to two charge $`1/3`$ quarks, $`D`$ and $`S`$. The masses of the two heavy quarks $`Q_{1,2}`$ are given, for either fermion embedding, by $$M_{Q_m}=s_\beta \lambda _{Q_m}f(m=1,2),$$ (15) where we have neglected the masses of the quarks of the first two generations and chosen $`\lambda _{Q_m}`$ to be the Yukawa coupling involving $`\mathrm{\Phi }_2`$ (see Sec. B.2.3 and B.2.6 for further details). The heavy quark couplings to the Higgs boson are proportional to the Yukawa couplings $`\lambda _{Q_m}`$ as expected, and can be rewritten in terms of the heavy quark mass $`M_Q`$ (see Table 2). After EWSB, the Yukawa couplings $`\lambda _{Q_m}`$ lead to mixing between the heavy quarks $`Q`$ and the corresponding SM quarks of like charge given by $`Q=Q_0\delta _qq_0`$, where as usual $`Q_0,q_0`$ denote the electroweak eigenstates of each generation. The mixing angle $`\delta _q`$ is given to order $`v/f`$ by $$\delta _q=\pm \frac{v}{\sqrt{2}ft_\beta }\delta _\nu ,$$ (16) where the upper sign is for the anomaly-free embedding ($`Q=D,S`$) and the lower sign is for the universal embedding ($`Q=U,C`$). The mixing between SM quarks and their heavy counterparts causes isospin violation at order $`\delta _\nu ^2`$ in processes involving only SM fermions. This isospin violation can be suppressed by choosing $`t_\beta 1`$. As in the top sector, the mixing due to $`\delta _q`$ gives rise to the couplings of $`Q`$ to $`q^{}W`$ and $`qZ`$; the Feynman rules are given in Table 2. Although the new heavy quarks $`Q_{1,2}`$ of the first two generations do not play a significant role in the cancellation of the Higgs mass quadratic divergence (they take part in the cancellation of the numerically insignificant Higgs mass quadratic divergence from their SM partners in the first two generations), they share the common parameters $`f`$ and $`t_\beta `$ with the top sector, providing additional experimental observables that can be used to test the little Higgs structure of the couplings. The new heavy quarks of the first two generations introduce two further parameters, which can be chosen as their masses $`M_{Q_m}`$ or equivalently their Yukawa couplings $`\lambda _{Q_m}`$, as related by Eq. (15). The couplings between the new heavy quarks and the TeV-scale gauge bosons are fixed by the gauge symmetry; they are summarized in Table 2. We will not comment on them further here since they will not play a significant role in our phenomenological analysis. ### 3.4 Heavy quark production and decay at the LHC #### 3.4.1 $`T`$ production and decay The top-partner $`T`$ can be pair-produced via QCD interactions at the LHC; however, because the final state contains two heavy particles, the pair-production cross section falls quickly with increasing $`M_T`$. Instead, single $`T`$ production via $`Wb`$ fusion yields a larger cross section in both the Littlest Higgs model and the SU(3) simple group model, as shown in Figs. 2 and 3, respectively. In the Littlest Higgs model, the single $`T`$ production cross section at fixed $`M_T`$ depends on only one model parameter, $`x_\lambda `$, as shown in Fig. 2. In particular, the cross section is proportional to $`x_\lambda ^2`$, as can be seen by examining the $`W^+\overline{T}b`$ coupling in Table 1 while holding $`M_T`$ fixed. We see that the cross section is typically in the range 0.01–100 fb for $`M_T=1.5`$–3.5 TeV. In the SU(3) simple group model, the single $`T`$ production cross section at fixed $`M_T`$ depends on two model parameters, $`x_\lambda `$ and $`t_\beta `$. From the $`W^+\overline{T}b`$ coupling in Table 1 one can see that at fixed $`M_T`$, the cross section scales with $`\lambda _T^2`$: $$\sigma \lambda _T^2s_\beta ^2c_\beta ^2(x_\lambda x_\lambda ^1)^2.$$ (17) The cross section is invariant under $`t_\beta 1/t_\beta `$ and under $`x_\lambda 1/x_\lambda `$. It reaches a maximum at $`t_\beta =1`$, and vanishes at $`x_\lambda =1`$. Away from unity, it falls like $`t_\beta ^2(t_\beta ^2)`$ for large (small) $`t_\beta `$, and grows like $`x_\lambda ^2(x_\lambda ^2)`$ for large (small) $`x_\lambda `$. The cross section is shown in Fig. 3 for $`t_\beta =3`$ and various values of $`x_\lambda `$. We see that the cross section is similar in size to that in the Littlest Higgs model, depending on the parameter values in either model. The dominant decay modes of $`T`$ in all little Higgs models are $`tH`$, $`tZ`$ and $`bW`$. The partial widths of $`T`$ to these final states are all controlled by the same coupling $`\lambda _T`$, $$\mathrm{\Gamma }(TtH)=\mathrm{\Gamma }(TtZ)=\frac{1}{2}\mathrm{\Gamma }(TbW)=\frac{\lambda _T^2}{32\pi }M_T=9.9\lambda _T^2\left(\frac{M_T}{\mathrm{TeV}}\right)\mathrm{GeV},$$ (18) where we neglect final-state masses compared to $`M_T`$. If these are the only decays of $`T`$, then its total width is $`40\lambda _T^2(M_T/\mathrm{TeV})`$ GeV. The branching fractions of $`T`$ into these final states are then given by $$\mathrm{BR}(TtH)=\mathrm{BR}(TtZ)=1/4,\mathrm{BR}(TbW)=1/2.$$ (19) This simple relation between the branching fractions is easily understood in terms of the Goldstone boson equivalence theorem: the decay modes at high energies (large $`M_T`$) are just those into the four components of the SM Higgs doublet, i.e., the three Goldstone degrees of freedom and the physical Higgs boson. Phenomenological studies of these $`T`$ decays have been performed at the level of somewhat realistic detector simulations in Ref. . The $`T`$ mass can be reconstructed from each of these three channels; $`TZt\mathrm{}^+\mathrm{}^{}b\mathrm{}\overline{)}E_T`$ provides the cleanest mass peak . If the only significant decays of $`T`$ are into $`tH`$, $`tZ`$ and $`bW`$, then the branching fractions of $`T`$ are predicted independent of any model parameters by Eq. (19). A measurement of the rate for single $`T`$ production with decays into any one of the three final states is sufficient to determine the production cross section, and thus extract $`\lambda _T`$. The measurement of the characteristic pattern of branching fractions also provides a test of the model (see Sec. 3.6.1). In the SU(3) simple group model, $`T`$ has additional possible decay modes due to the additional particles in the spectrum. In particular, $`T`$ can also decay to $`t\eta `$, $`tY^0`$, and $`bX^+`$ final states, depending on the relative masses of $`T`$, $`\eta `$, and $`X,Y`$. In order to measure the single $`T`$ production cross section, and hence $`\lambda _T`$, one needs to know the branching fraction(s) of the decay mode(s) in which $`T`$ is observed. Assuming the SU(3) simple group model structure, these can be predicted as follows. The $`T`$ mass can be reconstructed in, e.g., $`TZt\mathrm{}^+\mathrm{}^{}b\mathrm{}\overline{)}E_T`$ as discussed above. The $`X,Y`$ gauge boson masses are fixed in terms of $`M_Z^{}`$, which will be easily measurable from its decays to dileptons (see Sec. 4). The $`T`$ partial widths to $`tY`$ and $`bX`$ can then be calculated in terms of the gauge couplings in Table 2. The $`T`$ partial width to $`\eta `$ can be calculated from the coupling in Table 2 once the $`\eta `$ mass is measured, e.g., in decays of $`\eta `$ to dijets. The partial widths to $`tH`$, $`tZ`$ and $`bW`$ are proportional to $`\lambda _T^2`$; thus the only remaining free parameter to be extracted from the rate measurement in any given final state is $`\lambda _T`$. Measurements of the pattern of branching fractions then provide a nontrivial test of the model. Similarly, in the Littlest Higgs model with two U(1) groups gauged, $`T`$ can decay into $`tA_H`$. Once the $`A_H`$ mass is measured, a similar analysis can be applied. #### 3.4.2 $`Q`$ production and decay The heavy quarks $`Q`$ in the SU(3) simple group model can be produced at the LHC via, e.g., $`WdU`$, $`ZuU`$. The production couplings are given in Table 2; for fixed $`M_Q`$, the cross section depends on only one model parameter, $`\delta _\nu `$; in particular the cross section is proportional to $`\delta _\nu ^2=v^2/2f^2t_\beta ^2`$. The single production cross section for $`U+\overline{U}`$ is shown in Fig. 4, together with the $`U\overline{U}`$ pair production cross section from QCD. The single $`U`$ production cross section is quite large compared to single production of $`T`$ at a comparable mass because $`T`$ production requires a $`b`$ quark in the initial state, while $`U`$ production proceeds from a valence $`u`$ or $`d`$ quark. By measuring both $`M_U`$ and the single $`U`$ production cross section, as well as $`f`$ from measurements in the gauge sector (see Sec. 4), one can determine $`\lambda _U`$ and $`t_\beta `$ from Eqs. (15) and (16). This measurement of $`t_\beta `$ is independent from that in the $`T`$ sector and can be used as a nontrivial test of the model, as will be discussed further in Sec. 3.5. Production of the heavy quark partners of the first generation offers an additional powerful handle on the SU(3) simple group model. First, consider single $`U`$ production in the universal fermion embedding. This proceeds via the subprocesses $$dW^+U,uZU;\overline{d}W^{}\overline{U},\overline{u}Z\overline{U}.$$ (20) At a proton-proton collider such as the LHC, we expect the cross section for $`U`$ production, from initial-state valence $`u`$ and $`d`$ quarks, will be much larger than that for $`\overline{U}`$, from initial-state sea $`\overline{u}`$ and $`\overline{d}`$ antiquarks. In fact, $`\overline{U}`$ production constitutes less than 10% of the total $`U+\overline{U}`$ cross section shown in Fig. 4. There will thus be a large asymmetry in the charge of the final lepton in $`U,\overline{U}`$ decays to $`W^\pm `$, with many more positively charged leptons. In the anomaly-free embedding, single $`D`$ production proceeds via the subprocesses $$uW^{}D,dZD;\overline{u}W^+\overline{D},\overline{d}Z\overline{D}.$$ (21) Because of the parton densities in the proton, the rate for $`D`$ production via charged current will be somewhat higher than for $`U`$, while the rate for $`D`$ production via neutral current will be somewhat lower than for $`U`$, resulting in a comparable total cross section. Again, there will be a large asymmetry in the charge of the final lepton in $`D,\overline{D}`$ decays to $`W^{}`$, with many more negatively charged leptons. This allows a simple measurement of the dominant lepton charge in $`Qq^{}W(\mathrm{}\nu )`$ decays to distinguish the universal fermion embedding from the anomaly-free fermion embedding. The fermion embedding must be known in order for the model parameters to be extracted from the single-$`Q`$ production cross section because the embedding determines which parton densities enter the production cross section calculation. Just as for $`T`$, the decay modes of $`U`$ in the SU(3) simple group model depend on the spectrum of masses. The $`U`$ quark decays into $`uH`$, $`uZ`$ and $`dW`$ with partial widths $$\mathrm{\Gamma }(UuH)=\mathrm{\Gamma }(UuZ)=\frac{1}{2}\mathrm{\Gamma }(UdW)=5.0\left(\frac{\mathrm{TeV}}{ft_\beta }\right)^2\left(\frac{M_U}{\mathrm{TeV}}\right)^3\mathrm{GeV}.$$ (22) $`U`$ can also decay into $`u\eta `$; however, the coupling at leading order in $`v/f`$ is proportional to the up quark Yukawa coupling, so this decay is extremely suppressed and can be neglected. If $`U`$ is heavy enough, it can also decay into $`uY`$ and $`dX`$ with partial widths that depend only on the heavy gauge boson mass $`M_{X,Y}`$; the $`UuY`$ and $`UdX`$ couplings are fixed in terms of the SM gauge coupling $`g`$. The heavy gauge boson mass $`M_{X,Y}`$ can be obtained from the $`Z^{}`$ mass measurement (see Sec. 4). The partial widths to $`uH`$, $`uZ`$ and $`dW`$ can then be extracted together with $`\delta _\nu `$ from the rate measurement into any final state. The above discussion applies equally to $`D`$ in the anomaly-free fermion embedding. The signal kinematics are as follows. $`U`$ is produced via $`dW^+`$ or $`uZ`$ fusion, yielding a forward jet from which the $`W`$ or $`Z`$ was radiated. $`U`$ then decays into a high-$`p_T`$ quark and a $`W`$ boson, with $`W\mathrm{}\nu `$. The $`W`$ is highly boosted, with a momentum of roughly half the $`U`$ mass, so that the momenta of the neutrino and charged lepton are almost parallel. The decay kinematics are sketched in Fig. 5. We can take advantage of the large boost of the $`W`$ boson in $`U`$ decay to reconstruct the $`U`$ mass. Normally such a decay involving a neutrino in the final state would allow only the reconstruction of the $`U`$ transverse mass. However, because $`U`$ is very heavy, we can neglect the $`W`$ mass relative to its momentum and approximate the direction of the neutrino momentum to be parallel to that of the charged lepton. We can then reconstruct the full neutrino momentum and combine it with that of the charged lepton and the high-$`p_T`$ jet to reconstruct a mass peak for $`U`$. We apply the following cuts to select $`U`$ production events over the SM $`W^+jj`$ background. We require a positively-charged electron or muon with $$\|\eta _{\mathrm{}}\|<3,p_T\mathrm{}>20\mathrm{GeV}.$$ (23) For the central high-$`p_T`$ jet we require $$\|\eta _{j_1}\|<3,p_{Tj_1}>300\mathrm{GeV}.$$ (24) We also require that the forward jet be tagged, with $$3<\|\eta _{j_2}\|<5,p_{Tj_2}>30\mathrm{GeV}.$$ (25) Finally we require missing transverse momentum, $$p_T/>30\mathrm{GeV}.$$ (26) To simulate the detector effects, we smear the energies for the charged lepton and the jets according to a Gaussian form, $`\mathrm{\Delta }E/E=a/\sqrt{E/\mathrm{GeV}}b`$, with $`a=5\%`$, $`b=1\%`$ for a charged lepton and $`a=50\%`$, $`b=2\%`$ for a jet. The $`p_T`$ distribution of the highest-$`p_T`$ jet is shown in the left panel of Fig. 6, together with the $`W^+jj`$ background. The signal distribution clearly exhibits a Jacobian peak near $`M_U/2`$. The right panel of Fig. 6 shows the $`U`$ transverse mass and the fully reconstructed $`U`$ mass. The $`U`$ mass is reconstructed from the momenta of $`\mathrm{}^+`$ and the highest-$`p_T`$ jet, as well as the missing momentum assumed to point along the direction of the $`\mathrm{}^+`$ momentum. The reconstructed mass variable indeed leads to a sharper peak than the transverse mass. In Fig. 6 we have included only $`U`$ production (without the $`\overline{U}`$ contribution), and folded in the branching fractions of $`UW^+u`$ and $`W^+\mathrm{}^+\nu `$, with $`\mathrm{}^+=e^+,\mu ^+`$. The signal cross section after cuts for $`M_U=3`$ TeV and $`ft_\beta =3`$ TeV is about 0.66 fb, resulting in close to 200 signal events in 300 fb<sup>-1</sup> of LHC luminosity. The background is well under control. Additional statistics can be gained by considering the decay channels $`UuZ,uH`$. One can do a similar analysis for single $`C`$ ($`S`$) production, using $`M_C`$ ($`M_S`$) and the production cross section together with $`f`$ from the gauge sector measurements to determine $`\lambda _C`$ ($`\lambda _S`$) and make another independent measurement of $`t_\beta `$. However, because $`C`$ ($`S`$) is produced from inital-state sea quarks $`c`$ and $`s`$, its production rate will be lower, only 10–20% of that of $`U`$ ($`D`$). Further, since the sea quark and antiquark distributions are equal, there will be no asymmetry in the charge of the final lepton in $`C`$ ($`S`$) decays to $`W^\pm `$. This allows the $`C`$ ($`S`$) resonance to be experimentally distinguished from the $`U`$ ($`D`$) resonance, if enough events can be collected above background. ### 3.5 Testing the Higgs mass divergence cancellation in the top sector The key experimental test of the little Higgs models is to verify the cancellation of the Higgs mass quadratic divergence, embodied in the crucial relation of Eq. (13). Ideally, one could hope to measure the couplings $`\lambda _T`$ and $`\lambda _T^{}`$ directly, without making any assumptions about the model structure. The coupling $`\lambda _T`$ controls the $`T`$ production cross section in $`Wb`$ fusion, where it can be extracted by measuring the single-$`T`$ production rate and the $`T`$ mass from signal kinematics. The coupling $`\lambda _T^{}`$ could in principle be extracted from a measurement of the associated $`TH`$ production cross section. However, a quick estimate indicates that the cross section is too small to be observable at the LHC. Instead, the relation in Eq. (13) for the Higgs mass divergence cancellation must be checked within the context of the particular model. Once the model is determined, the relevant independent parameters that control the top sector must be overconstrained to make a nontrivial test of the model. In the Littlest Higgs model, one can use the model relation $`\lambda _T^{}=\lambda _TM_T/f`$ to write the divergence cancellation condition in terms of the four observables $`(\lambda _t,\lambda _T,M_T,f)`$. Note that only three of these are independent in the Littlest Higgs model; $`\lambda _T`$ and $`M_T`$ can both be written in terms of $`f`$, $`\lambda _t`$ and $`x_\lambda `$. Combining $`T`$-sector measurements of $`M_T`$ and $`\lambda _T`$ with a measurement of $`f`$ from the heavy gauge boson sector, one can overconstrain the parameters and verify the cancellation of the quadratic divergence. In the SU(3) simple group model the situation is more complicated because of the ratio of the two vacuum condensates, $`f_2/f_1=t_\beta `$, which appears in the fermion sector of the model. Thus, in addition to the four parameters $`(\lambda _t,\lambda _T,M_T,f)`$ measurable in the $`T`$ and heavy gauge boson sectors, one needs a measurement of $`t_\beta `$ in order to overconstrain the parameters and verify the relation in Eq. (13). Fortunately, $`t_\beta `$ can be extracted independently of the $`\lambda _T`$ and $`M_T`$ measurements by measuring the mass and production cross section of the $`U`$ or $`D`$ quarks, since their production couplings are proportional to $`1/t_\beta `$. ### 3.6 Comparison with other models #### 3.6.1 A fourth generation sequential top-prime The key feature that distinguishes $`T`$ from a fourth generation sequential top-prime is the fact that it is an SU(2) singlet before mixing with the top quark. This feature allows for the presence of a vectorlike mass term for $`T`$ and flavor-changing $`TtH`$ and $`TtZ`$ couplings in the mass basis, both of which are forbidden by electroweak symmetry in a fourth-generation model. As pointed out in Ref. , detecting and measuring the flavor-changing neutral current decays $`TZt`$ and $`THt`$, with equal branching fractions, allows one to rule out the fourth-generation hypothesis and conclude that $`T`$ is an electroweak singlet, acquiring its coupling to the Higgs via a gauge-invariant $`TtH`$ term. #### 3.6.2 The top quark see-saw In the top quark see-saw model , EWSB occurs via the condensation of the top quark in the presence of an extra vectorlike SU(2)-singlet quark, forming a composite Higgs boson. In order to reproduce the correct electroweak scale, the condensate mass must be large, of order 600 GeV. The vectorlike singlet quark joins the top in a see-saw, yielding the physical top mass (adjusted to the experimental value) and a multi-TeV mass for the vectorlike quark. The little Higgs models thus generically contain an extended top sector with the same electroweak quantum numbers as in the top see-saw model, i.e., a (multi-)TeV-scale isosinglet vectorlike quark $`T`$ with a small mixing with the SM top quark that gives rise to $`TtZ`$, $`TtH`$ and $`TbW`$ couplings. The most important difference between the top see-saw model and the little Higgs models is that the top see-saw model makes no prediction for the dimension-5 $`HHTT`$ coupling $`\lambda _T^{}`$, although this coupling can be generated radiatively. Thus, the top see-saw model does not in general satisfy the condition for cancellation of the Higgs mass quadratic divergence given in Eq. (13). In the top see-saw model, the $`TtH`$ coupling $`\lambda _T`$ is constrained by the compositeness condition, which requires the wavefunction renormalization of the composite Higgs field to vanish at the compositeness scale $`M_c`$. Ignoring the effect of EWSB, the effective Lagrangian of the top see-saw model is $$=Z_h\|𝒟h\|^2+[\sqrt{2}y_t\overline{\psi }_Lt_R\sqrt{Z_h}h+\sqrt{2}\lambda _T\overline{\psi }_LT_R\sqrt{Z_h}hM_T\overline{T}_LT_R+\mathrm{h}.\mathrm{c}.]+V_h,$$ (27) where $`Z_h`$ is the wavefunction renormalization of the composite Higgs field $`h`$ and $`V_h`$ is the usual SM Higgs potential. In the large-$`N_c`$ approximation, this implies $$\lambda _T^2=\frac{4\pi ^2}{N_c\mathrm{log}(M_c/M_T)}\frac{m_t^2}{v^2}.$$ (28) The compositeness scale $`M_c`$ should not be too far away from the scale of the heavy states. For $`M_c/M_T10`$–100 and $`N_c=3`$, we obtain $`\lambda _T5.2`$–2.4; in particular, the compositeness condition generally requires a fairly large value for $`\lambda _T`$. In little Higgs models, on the other hand, $`\lambda _T`$ is typically of order one or smaller. In the Littlest Higgs model, $`\lambda _T=x_\lambda m_t/vx_\lambda /\sqrt{2}`$, which reaches the typical top quark see-saw values only for $`x_\lambda 4`$. Large values of $`x_\lambda `$ in the Littlest Higgs model tend to push up the $`T`$ mass, leading to greater fine tuning in the electroweak scale. In the SU(3) simple group model, $`\lambda _T=s_\beta c_\beta (x_\lambda x_\lambda ^1)m_t/v`$, which is further suppressed by the $`s_\beta c_\beta 1/2`$ factor in front. ## 4 The gauge sector Little Higgs models extend the electroweak gauge group at the TeV scale. The structure of the extended electroweak gauge group determines crucial properties of the little Higgs model, which can be revealed by studying the new gauge bosons at the TeV scale. Therefore, we continue with a study of the heavy gauge boson sectors of little Higgs models. ### 4.1 Heavy gauge boson masses and parameters The extra gauge bosons get their masses from the $`f`$ condensate, which breaks the extended gauge symmetry. For our two prototype models, the gauge boson masses are given in terms of the model parameters by $`\begin{array}{c}M_{W_H}=M_{Z_H}=gf/2sc=0.65f/\mathrm{sin}2\theta \hfill \\ M_{A_H}=gs_Wf/2\sqrt{5}c_Ws^{}c^{}=0.16f/\mathrm{sin}2\theta ^{}\hfill \end{array}\}\mathrm{in}\mathrm{the}\mathrm{Littlest}\mathrm{Higgs}\mathrm{model},`$ (31) $`\begin{array}{c}M_Z^{}=\sqrt{2}gf/\sqrt{3t_W^2}=0.56f\hfill \\ M_X=M_Y=gf/\sqrt{2}=0.46f=0.82M_Z^{}\hfill \end{array}\}\mathrm{in}\mathrm{the}\mathrm{SU}(3)\mathrm{simple}\mathrm{group}\mathrm{model}.`$ (34) In the SU(3) simple group model the heavy gauge boson masses are determined by only one free parameter, the scale $`f=\sqrt{f_1^2+f_2^2}`$. The Littlest Higgs model has two additional gauge sector parameters, $`\mathrm{tan}\theta =s/c=g_2/g_1`$ \[in the SU(2)$`{}_{}{}^{2}`$SU(2) breaking sector\] and $`\mathrm{tan}\theta ^{}=s^{}/c^{}=g_2^{}/g_1^{}`$ \[in the U(1)$`{}_{}{}^{2}`$U(1) breaking sector\]. If only one copy of U(1) is gauged , the $`A_H`$ state is not present and the gauge sector of the Littlest Higgs model is controlled by only two free parameters, $`f`$ and $`\mathrm{tan}\theta `$. Because the model with only one copy of U(1) gauged is favored by the electroweak precision constraints, and since the U(1) sectors of the product group models are quite model-dependent, we focus in what follows on the heavy SU(2) gauge bosons $`W_H`$ and $`Z_H`$. The $`W_H`$ and $`Z_H`$ bosons capture the crucial features of the gauge sector of the Littlest Higgs model and their phenomenology can be applied directly to the other product group models. ### 4.2 Heavy gauge boson interactions with SM particles The gauge couplings of the Higgs doublet take the general form $$=\{\begin{array}{c}\left[G_{HHVV}VV+G_{HHV^{}V^{}}V^{}V^{}+G_{HHVV^{}}VV^{}\right]H^2\hfill \\ \left[G_{HHV^+V^{}}V^+V^{}+G_{HHV^+V^{}}V^+V^{}+G_{HHV^+V^{}}(V^+V^{}+V^{}V^+)\right]H^2,\hfill \end{array}$$ (35) where the top line is for $`V`$ neutral and the bottom line is for $`V`$ charged. Here $`V`$ and $`V^{}`$ stand for the SM and heavy gauge bosons, respectively. This Lagrangian leads to two quadratically divergent diagrams contributing to the Higgs mass: one involving a loop of $`V`$, proportional to $`G_{HHVV}`$, and the other involving a loop of $`V^{}`$, proportional to $`G_{HHV^{}V^{}}`$. The divergence cancellation in the gauge sector can thus be written as $$\underset{i}{}G_{HHV_iV_i}=0,$$ (36) where the sum runs over all gauge bosons in the model. The couplings in the models under consideration are given in Table 3. In the SU(3) simple group model, the quadratic divergence cancels between the $`Z`$ and $`Z^{}`$ loops and between the $`W`$ and $`X`$ loops. In the Littlest Higgs model, the quadratic divergence cancels between the $`W`$ and $`W_H`$ loops and there is a partial cancellation between the $`Z`$ and $`Z_H`$ loops. Including the $`A_H`$ loop leads to a complete cancellation of the quadratic divergence from the $`Z`$ loop. The key test of the little Higgs mechanism in the gauge sector is the experimental verification of Eq. (36); we discuss the prospects further in Sec. 4.4. After EWSB, the couplings of $`H^2`$ to one heavy and one SM gauge boson induce mixing between the heavy and SM gauge bosons: $$V^{}=V_0^{}\delta _VV_0,\delta _V=v^2G_{HHVV^{}}/M_V^{}^2,$$ (37) where $`V_0^{}`$, $`V_0`$ stand for the states before EWSB. The mixing parameters $`\delta _V`$ are given in Table 3. This mixing gives rise to triple gauge couplings between one heavy and two SM gauge bosons, also shown in Table 3. In the SU(3) simple group model, EWSB also splits the $`X`$ and $`Y`$ gauge boson masses by a small amount, $$M_YM_X=\frac{gv^2}{4\sqrt{2}f}3.9\left(\frac{\mathrm{TeV}}{M_Z^{}}\right)\mathrm{GeV}.$$ (38) In the Littlest Higgs model, the couplings of the heavy gauge bosons to the SU(2)<sub>L</sub> fermion currents take the form $$Z_H^\mu \overline{f}f:ig\mathrm{cot}\theta T_f^3\gamma ^\mu P_L,W_H^{+\mu }\overline{u}d:\frac{ig}{\sqrt{2}}\mathrm{cot}\theta \gamma ^\mu P_L,$$ (39) where $`T_f^3=1/2`$ $`(1/2)`$ for up (down) type fermions. Below the TeV scale, exchange of $`W_H`$ and $`Z_H`$ gives rise to four-fermi operators, which are constrained by the electroweak precision data. The experimental constraints are loosened by going to small values of $`\mathrm{cot}\theta `$, for which the couplings of the heavy gauge bosons are suppressed. In the SU(3) simple group model, the $`Z^{}`$ couples to SM fermions with gauge strength, while the $`X,Y`$ gauge bosons couple only via the mixing between SM fermions and their TeV-scale partners. The couplings are given in Table 4. ### 4.3 Heavy gauge boson production and decay The best way to discover new heavy gauge bosons at the LHC is generally through Drell-Yan production. This is certainly true in the little Higgs models. In the Littlest Higgs model, the heavy gauge bosons $`Z_H,W_H`$ couple to pairs of SM fermions through the SU(2)<sub>L</sub> current, with coupling strength scaled by $`\mathrm{cot}\theta `$ compared to the SM SU(2)<sub>L</sub> couplings. They thus have large production cross sections, as shown in Fig. 7, controlled by one common free parameter, $`\mathrm{cot}\theta `$.<sup>5</sup><sup>5</sup>5Note that the electroweak precision data tend to favor small values of $`\mathrm{cot}\theta `$, which reduces the contribution of $`W_H,Z_H`$ to four-Fermi operators at low energy. Small $`\mathrm{cot}\theta `$ lowers the Drell-Yan cross section, reducing the LHC reach for $`W_H,Z_H`$ discovery. In addition, because $`Z_H`$ and $`W_H`$ form an SU(2) triplet, they are degenerate in mass up to very small EWSB effects. Thus, the measurement of the $`Z_H`$ mass in dileptons predicts the transverse mass distribution of the $`W_H`$ in $`W_H\mathrm{}\nu `$, and the measurement of the rate for $`Z_H`$ into dileptons predicts the rate for $`W_H`$ into leptons, allowing a test of the SU(2) triplet nautre of $`(W_H,Z_H)`$. In the SU(3) simple group model, the heavy gauge boson $`Z^{}`$ couples to pairs of SM fermions with couplings fixed in terms of the SM gauge couplings and depending only on the (discrete) choice of the fermion embedding, as shown in the left panel of Fig. 7. Unlike the $`Z_H`$ of the Littlest Higgs model, there is no tunable parameter in the $`Z^{}`$ cross section.<sup>6</sup><sup>6</sup>6This parameter independence is the most characteristic feature of the $`Z^{}`$ in simple group models with the extended gauge group SU(3)$`\times `$U(1)<sub>X</sub> . Models with a larger extended gauge group, SU($`N`$)$`\times `$U(1)<sub>X</sub> with $`N>3`$, lose this parameter independence because they contain $`N2`$ broken diagonal generators, which mix in general. For example, the SU(4)$`\times `$U(1)<sub>X</sub> model of Ref. contains two broken diagonal generators, $`Z_1^{}`$ (which couples to SM fermion pairs with fixed strength) and $`Z_2^{}`$ (which does not couple to fermion pairs). After mixing, the mass eigenstates $`Z^{},Z^{\prime \prime }`$ share the fermion couplings with the mixing angle as a free parameter. If the fermion couplings of both states can be measured, the parameter independence reappears in the form of a coupling sum rule. The heavy gauge bosons $`X,Y`$ of the SU(3) simple group model have a very different phenomenology, rooted in their identity as the SU(2)<sub>L</sub> doublet $`(X^{},Y^0)`$ of broken off-diagonal generators of SU(3). Because they couple to SM quark pairs only through $`qQ`$ mixing as given in Table 4, their production cross sections in Drell-Yan are suppressed by $`\delta _\nu ^2v^2/f^2`$. This is shown for $`X`$ in the right panel of Fig. 7. Because of this large cross section difference, $`X^\pm `$ cannot be mistaken for the charged members of an SU(2) triplet containing $`Z^{}`$, providing an easy distinction between simple group and product group models. The $`20\%`$ mass splitting between $`X^\pm `$ and $`Z^{}`$ given in Eq. (34) also serves to distinguish $`X^\pm ,Z^{}`$ from an SU(2) triplet. An important feature of the product group models is the couplings of $`Z_H`$, $`W_H`$ to dibosons, which gives rise to the decays $`Z_HZH`$, $`W^+W^{}`$ and $`W_HWH`$, $`WZ`$. These couplings arise from a $`W_H^aW^ahh^{}`$ term in the Lagrangian and are proportional to $`\mathrm{cot}2\theta `$ due to the characteristic collective breaking structure of the gauge couplings in the product group models. The bosonic decay modes are dominated by the longitudinal components of the final-state bosons; their partial widths can be shown by the Goldstone boson equivalence theorem to obey the relation $`\mathrm{\Gamma }(Z_HZH)=\mathrm{\Gamma }(Z_HW^+W^{})=\mathrm{\Gamma }(W_HWH)=\mathrm{\Gamma }(W_HWZ)\mathrm{\Gamma }(V_HVH)`$, where we negect final-state masses and $$\mathrm{\Gamma }(V_HVH)=\frac{g^2\mathrm{cot}^22\theta }{192\pi }M_{V_H}=0.70\mathrm{cot}^22\theta \left(\frac{M_{V_H}}{\mathrm{TeV}}\right)\mathrm{GeV}.$$ (40) Here $`M_{V_H}`$ is the mass of $`Z_H`$ or $`W_H`$. The measurement of $`\mathrm{cot}\theta `$ from $`Z_H\mathrm{}^+\mathrm{}^{}`$ thus predicts the rates for decays of both $`Z_H`$ and $`W_H`$ into dibosons. The decay branching fractions of $`Z_H`$ and $`W_H`$ in the Littlest Higgs model are shown as a function of $`\mathrm{cot}\theta `$ in Fig. 8. We neglect final-state masses and assume that no decays to $`A_H`$ are present (namely, $`Z_HA_HH`$ and $`W_HA_HW`$). In the SU(3) simple group model, the decay partial widths of $`Z^{}`$ into pairs of SM bosons, $`ZH`$ and $`W^+W^{}`$, are fixed in terms of the $`Z^{}`$ mass (neglecting final-state masses) to be $$\mathrm{\Gamma }(Z^{}ZH)=\mathrm{\Gamma }(Z^{}W^+W^{})=\frac{g^2(1t_W^2)^2}{192\pi (3t_W^2)}M_Z^{}=0.13\left(\frac{M_Z^{}}{\mathrm{TeV}}\right)\mathrm{GeV},$$ (41) and the decay partial widths into pairs of SM fermions are fixed once the fermion embedding is chosen. As discussed in Sec. 3.4.2, the fermion embedding can be determined at the LHC by detecting the TeV-scale quark partner of the first generation, $`U`$ or $`D`$, decaying into $`Wq`$; the charge asymmetry of the final-state $`W`$ then determines the embedding. Knowledge of the fermion embedding from the fermion sector can be used to compute the $`Z^{}`$ couplings uniquely and perform a cross-check the model. If the TeV-scale fermion partners $`T`$ and/or $`Q_m`$ are not too heavy, they can be present in $`Z^{}`$ boson decays. If kinematically accessible, decays of $`Z^{}`$ to pairs of TeV-scale fermion partners proceed via gauge couplings. This is in contrast to the product group models, in which the TeV-scale top quark partner is mostly electroweak singlet and couples to $`Z_H`$ only through its electroweak doublet admixture at order $`v^2/f^2`$. The $`Z^{}`$ can also decay to one SM fermion and one TeV-scale fermion partner; however, the partial widths of these decays are suppressed by $`\delta _t^2,\delta _\nu ^2v^2/f^2`$ and will be numerically unimportant. Finally, the decay $`Z^{}Y^0\eta `$ will be kinematically accessible if $`\eta `$ is lighter than the $`Z^{}`$$`Y^0`$ mass splitting, $$M_Z^{}M_Y=0.18M_Z^{}=180\left(\frac{M_Z^{}}{\mathrm{TeV}}\right)\mathrm{GeV}.$$ (42) The decay branching fractions of $`Z^{}`$ in the SU(3) simple group model are given in Table 5, assuming that decays to TeV-scale fermion-partner pairs or to $`Y^0\eta `$ are kinematically forbidden and neglecting final-state masses. ### 4.4 Testing the Higgs mass divergence cancellation in the gauge sector The defining feature of the little Higgs models is the cancellation of the Higgs mass quadratic divergence at one-loop level. Here we investigate this cancellation in the gauge sector, as embodied in Eq. (36). Ideally, one could hope to measure directly the couplings $`G_{HHV^{}V^{}}`$ for each heavy gauge boson $`V^{}`$ in the model. This could be done by measuring associated production of $`H`$ with a heavy gauge boson; e.g., $`Z^{}H`$ associated production in the SU(3) simple group model. This probes $`G_{HHZ^{}Z^{}}`$ through the diagram involving $`q\overline{q}Z^{}Z^{}H`$, where one Higgs boson has been replaced by its vev in the interaction vertex. Ideally, one will want to measure both the magnitude and the sign of $`G_{HHZ^{}Z^{}}`$, perhaps through its interference with the similar diagram containing an $`s`$-channel $`Z`$. A detailed study is needed. In addition to testing the divergence cancellation, the measurement of the $`HHV^{}V^{}`$ couplings also sheds light onto the structure of the model by revealing which heavy gauge bosons are involved in the cancellation of each SM contribution to the Higgs mass quadratic divergence. In the Littlest Higgs model, $`Z_H`$ cancels the divergence from the SM $`W^3`$ boson, $`W_H^+`$ and $`W_H^{}`$ cancel the divergence from the SM $`W^\pm `$ bosons, and $`A_H`$ (if it is present) cancels the divergence from the SM hypercharge boson. In contrast, in the SU(3) simple group model, $`Z^{}`$ cancels the divergences from the SM $`W^3`$ boson and the hypercharge boson, while $`X`$ (together with its isospin partner $`Y`$) cancels the divergence from the SM $`W^\pm `$ bosons. Thus the $`HHZ^{}Z^{}`$ coupling strength that is characteristic of the little Higgs divergence cancellation mechanism can vary from model to model. In all product group models with SU(2)$`{}_{}{}^{2}`$SU(2)<sub>L</sub> breaking structure, the value of this coupling will be the same as in the Littlest Higgs model. In simple group models the value of the coupling will be different, and may depend on the model. For example, in the SU(4)$`\times `$U(1)<sub>X</sub> model of Ref. , the two broken diagonal generators mix to form mass eigenstates $`Z^{}`$ and $`Z^{\prime \prime }`$, which both take part in the divergence cancellation; the sum rule then reads $$G_{HHZZ}+G_{HHZ^{}Z^{}}+G_{HHZ^{\prime \prime }Z^{\prime \prime }}=0.$$ (43) A second approach to test the Higgs mass divergence cancellation, first described in Ref. , is to measure the couplings of Higgs bosons to one SM gauge boson and one new heavy gauge boson: e.g., $`HHW^+W_H^{}`$, $`HHZZ_H`$ in the Littlest Higgs model . This approach works only for the product group models, in which these couplings show a characteristic $`\mathrm{cot}2\theta `$ dependence which is fixed by the collective breaking structure of the gauge couplings and the nonlinear transformation of the SM Higgs doublet under the enlarged gauge symmetry. A “Big Higgs” model, in which the Higgs doublet transformed linearly under one of the two SU(2) gauge groups as the fermion doublets do, would have a $`HHZZ_H`$ coupling proportional to $`g\mathrm{cot}\theta `$ \[if $`h`$ transformed under SU(2)<sub>1</sub>\] or $`g\mathrm{tan}\theta `$ \[if $`h`$ transformed under SU(2)<sub>2</sub>\]. These couplings can be probed in the decays $`Z_HZH`$ and $`W_HWH`$ from $`Z_H,W_H`$ bosons produced on-shell, and will thus be more straightforward to measure than the $`HHV^{}V^{}`$ couplings discussed above. The $`\mathrm{cot}\theta `$ dependence of the $`Z_H`$ production cross section and decay to dileptons and the $`\mathrm{cot}2\theta `$ dependence of the $`Z_H`$ decay to $`ZH`$ can be probed simultaneously by measuring the rate into dileptons and the rate into $`ZH`$ ; these rates will fall upon the curve shown in Fig. 9 for the Littlest Higgs model. In simple group models, the $`HHZZ^{}`$ coupling does not provide a probe of the Higgs mass divergence cancellation because in these models this coupling is not directly related to the crucial $`HHZ^{}Z^{}`$ vertex that takes part in the cancellation of the Higgs mass quadratic divergence in the gauge sector. In fact, in the SU(3) simple group model, the $`HHZZ^{}`$ coupling is fixed by the extended gauge structure and would be the same in any model with the gauge group SU(3)$`\times `$U(1), whether or not the little Higgs mechanism were realized. The rates of $`Z^{}`$ into dileptons and into $`ZH`$ in the SU(3) simple group model are predicted uniquely for the universal and anomaly-free fermion embeddings, as shown in Fig. 9. In order to test the cancellation of the quadratic divergence in simple group models, it is thus very important to uncover the gauge structure and fermion embedding of the model. For this purpose, we now turn to a discussion of the determination of the $`Z^{}`$ properties in the simple group models. ### 4.5 Identifying the $`Z^{}`$ In addition to testing the little Higgs mechanism in the gauge sector as described in the previous section, one must also identify the model to which a newly-discovered $`Z^{}`$ boson belongs. This entails identifying the extended gauge structure and determining how the SM fits into it. We examine here some techniques that can be used at the LHC to shed light on the couplings of the $`Z^{}`$. We consider the $`Z_H`$ of the Littlest Higgs model and the $`Z^{}`$ of the SU(3) simple group model, with both the universal and anomaly-free fermion embeddings. As examples of other new physics possibilities, we also consider a sequential $`Z^{}`$ with the same couplings to fermions as the SM $`Z`$ boson, the $`Z_\psi ^{}`$ and $`Z_\chi ^{}`$ bosons of the $`E_6`$ model , and $`Z_R`$ of the left-right symmetric model . #### 4.5.1 Rate in dileptons A $`Z^{}`$ boson will most likely be first discovered in decays to dileptons. The dilepton rate then immediately tells us the production cross section times the leptonic branching ratio, and thus fixes a combination of the $`Z^{}`$ couplings to up and down quarks (in the production cross section), the $`Z^{}`$ coupling to leptons (in the decay partial width), and the $`Z^{}`$ total width (which enters the branching ratio to leptons). While the $`Z^{}`$ couplings to up and down quarks enter the production cross section together, multiplied by the appropriate parton densities, it may be possible to separate them experimentally by fitting the shape of the $`Z^{}`$ rapidity distribution to high-precision measurements of the up and down quark parton densities . The SU(3) simple group model gives a definite prediction for the $`Z^{}\mathrm{}^+\mathrm{}^{}`$ rate in each of the fermion embeddings, shown on the horizontal axis of Fig. 9. If extra decay modes of $`Z^{}`$ to the heavy fermion partners are kinematically allowed, they will increase the $`Z^{}`$ total width and thus decrease the rate into dileptons. Decays of $`Z^{}`$ into one SM and one heavy fermion are suppressed by the heavy-light mixing, $`v^2/f^2`$. Thus only decays into pairs of heavy fermions can contribute significantly; these are likely to be either kinematically inaccessible or heavily suppressed by phase space. In the Littlest Higgs model, the rate of $`Z_H`$ into dileptons depends on the free parameter $`\mathrm{cot}\theta `$. Thus, in this channel, the Littlest Higgs model can fake any other $`Z^{}`$ model for an appropriate value of $`\mathrm{cot}\theta `$. The rate in dileptons is uniquely predicted for the left-right symmetric model $`Z_R`$ and for a sequential $`Z^{}`$ (unless a tunable coupling is introduced by hand). The $`Z^{}`$ bosons in the $`E_6`$ model can mix, introducing a free parameter in their cross sections; however, the cross section is still constrained within a particular range for a $`Z^{}`$ of given mass, and the mixing angle can be extracted from the cross section. A $`Z^{}`$ from an extra U(1) gives a rate in dileptons tunable with the U(1) coupling. Therefore, while this rate measurement gives some valuable information about the $`Z^{}`$ couplings, it cannot uniquely determine the model. #### 4.5.2 Decay branching fractions to other fermion species In order to probe the $`Z^{}`$ couplings to fermions in more detail, one must look for $`Z^{}`$ decays into additional fermion species. This opens a window onto the relative couplings of the $`Z^{}`$ to particles with different hypercharges. Decays into neutrinos are only accessible through the $`Z^{}`$ total width, which in little Higgs models is typically smaller than the detector dilepton mass resolution (see Table 5). We thus consider decays into pairs of quarks. This is a more difficult search than detecting the $`Z^{}`$ in dileptons because of the large dijet background at the LHC. However, it may be possible to detect the $`Z^{}`$ decaying into top quark pairs, as a peak in the $`t\overline{t}`$ invariant mass spectrum, or into bottom quark pairs, as a peak in the $`b`$-tagged dijet invariant mass spectrum. Measuring the rate of the $`Z^{}`$ into top (bottom) quark pairs and taking the ratio with the rate to dileptons gives the ratio of partial widths into top (bottom) versus electrons, shown in Table 6. In the Littlest Higgs model, this ratio is fixed independent of $`\mathrm{cot}\theta `$ because the $`\mathrm{cot}\theta `$ dependence enters the couplings to all fermions in the same way. Further, because $`Z_H`$ couples universally to all fermion doublets, this ratio is just given by the number of color degrees of freedom, $`N_c=3`$ (neglecting final-state masses). This ratio is also fixed in the SU(3) simple group model; it is different from the value in the Littlest Higgs model because of the U(1)<sub>X</sub> content of the $`Z^{}`$, which introduces a dependence on the fermion hypercharge. Note that the ratio of top (bottom) to electron partial widths is the same in the universal and the anomaly-free fermion embeddings, because in both embeddings the leptons and the third generation of quarks all transform as $`\mathrm{𝟑}`$s of SU(3); the difference between the two embeddings appears only in the first two generations of quarks. Similarly, these ratios are independent of model parameters for a sequential $`Z^{}`$, the $`E_6`$ $`Z_\psi ^{}`$ and $`Z_\chi ^{}`$, and the left-right symmetric $`Z_R`$. The $`E_6`$ $`Z_\psi ^{}`$ and $`Z_\chi ^{}`$ mix in general, leading to intermediate values of the partial width raitos. $`Z_\psi ^{}`$ has the same $`\mathrm{BR}(tt)/\mathrm{BR}(ee)`$ and $`\mathrm{BR}(bb)/\mathrm{BR}(ee)`$ as the Littlest Higgs $`Z_H`$, and $`Z_\chi ^{}`$ has the same $`\mathrm{BR}(bb)/\mathrm{BR}(ee)`$, as the Littlest Higgs $`Z_H`$. Likewise, the sequential $`Z^{}`$ has the same $`\mathrm{BR}(bb)/\mathrm{BR}(ee)`$ as the SU(3) simple group model $`Z^{}`$; however, its $`\mathrm{BR}(tt)/\mathrm{BR}(ee)`$ is rather different. Of course, the couplings of a $`Z^{}`$ from an anomalous extra U(1) can be tuned to duplicate the predictions of any of these models. #### 4.5.3 Forward-backward asymmetry The forward-backward asymmetry in $`f_i\overline{f}_iZ^{}f_f\overline{f}_f`$ probes the chiral structure of the $`Z^{}`$ couplings to the initial- and final-state fermions. At the partonic level, this asymmetry is defined as $$A_{FB}^{0,if}=\frac{N_FN_B}{N_F+N_B}=\frac{3}{4}𝒜_i𝒜_f,$$ (44) where $`N_F`$ ($`N_B`$) is the number of events with the final-state fermion momentum in the forward (backward) direction defined relative to the initial-state fermion. The asymmetry $`𝒜_f`$ is defined in terms of the couplings $`g_{L,R}^f`$ as $$𝒜_f=\frac{(g_L^f)^2(g_R^f)^2}{(g_L^f)^2+(g_R^f)^2}.$$ (45) Even though the LHC is a symmetric $`pp`$ collider, a forward-backward asymmetry can be defined by taking advantage of the fact that the valence quarks in the proton tend to carry a higher momentum fraction $`x`$ than the sea (anti)quarks . A “hadronic” forward-backward asymmetry can then be defined as $$A_{FB}^{\mathrm{had}}=\frac{N_FN_B}{N_F+N_B},$$ (46) where now the forward direction for the final-state fermion is defined relative to the boost direction of the $`Z^{}`$ center-of-mass frame. In the narrow-width approximation (neglecting interference between the $`Z^{}`$ resonance and the continuum photon and $`Z`$ exchange), $`A_{FB}^{\mathrm{had}}`$ is given in terms of the partonic asymmetries by $$A_{FB}^{\mathrm{had}}=\frac{𝑑x_1_{q=u,d}A_{FB}^{0,qf}\left(F_q(x_1)F_{\overline{q}}(x_2)F_{\overline{q}}(x_1)F_q(x_2)\right)\mathrm{sign}(x_1x_2)}{𝑑x_1_{q=u,d,s,c}\left(F_q(x_1)F_{\overline{q}}(x_2)+F_{\overline{q}}(x_1)F_q(x_2)\right)},$$ (47) where $`F_q(x_1)`$ is the parton distribution function (PDF) for quark $`q`$ in the proton with momentum fraction $`x_1`$, evaluated at $`Q^2=M_Z^{}^2`$. The momentum fraction $`x_2`$ is related to $`x_1`$ by the condition $`x_1x_2=M_Z^{}^2/s`$ in the narrow-width approximation. Only $`u`$ and $`d`$ quarks contribute to the numerator since we explicitly take the quark and antiquark PDFs to be identical for the sea quarks; all flavors contribute to the denominator. Here we consider $`Z^{}`$ decays to $`e^+e^{}`$ only, since it is much easier at LHC to determine the charge of a lepton than the charge of a quark. Decays to $`\mu ^+\mu ^{}`$ can be added to double the statistics. The relevant partonic asymmetries and $`A_{FB}^{\mathrm{had}}`$ are listed in Table 7 for the little Higgs models under consideration, as well as a number of other $`Z^{}`$ models. The hadronic forward-backward asymmetry $`A_{FB}^{\mathrm{had}}`$ varies with $`M_Z^{}`$ due to the shape of the PDFs. The $`Z^{}`$ mass dependence is shown in Fig. 10 for the models included in Table 7. It is interesting to note that the asymmetries of the $`E_6`$ $`Z^{}`$ bosons are less than or equal to zero, unlike the rest of the models. The $`E_6`$ boson asymmetries remain negative definite for arbitrary mixing between $`Z_\psi ^{}`$ and $`Z_\chi ^{}`$: $`A_{FB}^{0,ue}`$ is always zero and $`A_{FB}^{0,de}`$ varies between $`0.75`$ and 0 depending on the mixing angle. In Eq. (47) we have expressed $`A_{FB}^{\mathrm{had}}`$ as a single number, integrated over rapidity, which depends on both $`A_{FB}^{0,ue}`$ and $`A_{FB}^{0,de}`$. It may be possible to extract these two quantities separately by fitting the asymmetry as a function of the $`Z^{}`$ rapidity to high-precision measurements of the up and down quark parton densities ; however, this would require a huge amount of luminosity. In the Littlest Higgs model, a measurement of $`A_{FB}^{\mathrm{had}}`$ would provide a spectacular test of the model because it would confirm that $`𝒜_u=𝒜_d=𝒜_e=\pm 1`$; that is, that the $`Z_H`$ couplings to fermions are either purely left-handed or purely right-handed. The sign ambiguity is due to the fact that $`A_{FB}^{0,if}`$ depends on the product $`𝒜_i𝒜_f`$. Together with measurements of $`\mathrm{BR}(tt)/\mathrm{BR}(ee)`$ and/or $`\mathrm{BR}(bb)/\mathrm{BR}(ee)`$, which would demonstrate the universality of the $`Z_H`$ couplings to fermions, and the discovery of the $`W_H^\pm `$ degenerate in mass and with a related production rate, this measurement would confirm $`Z_H`$ as a member of an SU(2) triplet of gauge bosons. In such a case we learn that the SM SU(2)<sub>L</sub> gauge symmetry arises from the diagonal breaking of \[SU(2)\]<sup>2</sup>, with the SM fermion doublets transforming under one of the two SU(2) gauge groups. A measurement of $`A_{FB}^{\mathrm{had}}`$ will also provide a test of the SU(3) simple group model and the other $`Z^{}`$ models considered, since it probes another independent combination of the $`Z^{}`$ couplings to fermions. #### 4.5.4 Bosonic decay modes Measuring the bosonic decay modes $`Z^{}ZH`$ and $`Z^{}W^+W^{}`$ probes the transformation properties of the Higgs doublet under the extended gauge symmetry and the mixing of $`Z`$ and $`Z^{}`$ induced by electroweak symmetry breaking. As described in detail in Sec. 4.4, this can shed light on the little Higgs mechanism in the gauge sector, but it also provides useful information about the model structure. Also of interest are bosonic decay modes of the $`Z^{}`$ involving non-SM bosons in the final state, such as $`Z^{}Y\eta `$ in the SU(3) simple group model or $`Z_HA_HH`$ in the Littlest Higgs model. Detecting and measuring the branching fractions of these decay modes provides additional information on the structure of the extended gauge group and the mixings among the new gauge bosons. ## 5 Other phenomenological features of the SU(3) simple group model In this section we collect some additional features of the SU(3) simple group model not directly relevant to the simple group/product group classification and the identification of the little Higgs mechanism. ### 5.1 The heavy leptons In the SU(3) simple group model, the three lepton doublets of the SM are enlarged into triplets. The model thus contains three heavy neutral states $`N_m`$. The scalar interactions of the leptons can be written as $$_Y=i\lambda _{N_m}N_m^c\mathrm{\Phi }_2^{}L_m+\frac{i\lambda _e^{mn}}{\mathrm{\Lambda }}e_m^cϵ_{ijk}\mathrm{\Phi }_1^i\mathrm{\Phi }_2^jL_n^k+\mathrm{h}.\mathrm{c}.,$$ (48) where $`m,n=1,2,3`$ are generation indices, $`i,j,k=1,2,3`$ are SU(3) indices, $`L_m=(\nu ,e,iN)_m^T`$ are the lepton triplets, and $`N_m^c`$ are right-handed neutral leptons that marry the $`N_m`$ and get masses of order $`f`$ TeV. We neglect neutrino masses; a nice extension of the SU(3) simple group model including neutrino masses was presented in Ref. . Equation (48) generates masses for $`N_m`$, $$M_{N_m}=\lambda _{N_m}s_\beta f.$$ (49) The Lagrangian also contains a term $$_Y\frac{\lambda _{N_m}c_\beta }{\sqrt{2}}HN_m^c\nu +\mathrm{h}.\mathrm{c}.=\frac{M_{N_m}}{\sqrt{2}ft_\beta }HN_m^c\nu +\mathrm{h}.\mathrm{c}.$$ (50) for each generation, leading to mixing between the $`N_m`$ and the SM neutrinos given by $`N=N_0\delta _\nu \nu _0`$, where $`N_0,\nu _0`$ denote the electroweak eigenstates of each generation and $`\delta _\nu `$ was given in Eq. (16). This mixing gives rise to the couplings of $`N`$ to $`eW`$ and $`\nu Z`$ with Feynman rules $$W_\mu ^+\overline{N}e:\frac{ig\delta _\nu }{\sqrt{2}}\gamma _\mu P_L,Z_\mu \overline{N}\nu :\frac{ig\delta _\nu }{2c_W}\gamma _\mu P_L.$$ (51) Because the $`N_m`$ carry lepton number, their production at the LHC requires an additional lepton in the final state and can thus proceed only through $`s`$-channel gauge boson exchange, e.g., $`q\overline{q}^{}W^+Ne^+`$. Their decays, into $`\nu H`$, $`eW`$ and $`\nu Z`$, along with $`eX`$, $`\nu Y`$ and $`\nu \eta `$ if kinematically accessible, will be spectacular. The $`N_m`$ could also be produced at a linear collider of sufficient energy through $`t`$-channel $`W`$ exchange, $`e^+e^{}\overline{\nu }N`$. ### 5.2 The $`X`$ and $`Y`$ gauge bosons The heavy gauge bosons $`X^{},Y^0`$ correspond to the off-diagonal broken generators of SU(3) and thus communicate between the SU(2)<sub>L</sub> doublet fermions and the SU(2)<sub>L</sub> singlets, with couplings of gauge strength of the form $`XQq^{}`$ and $`YQq`$ as summarized in Table 2. These couplings can play a role in $`T`$ or $`Q`$ decay if the corresponding final states are kinematically accessible. They will not play a significant role in single $`T`$ or $`Q`$ production because the initial-state couplings of $`X^{},Y^0`$ to pairs of SM fermions are suppressed by $`v/f`$. While $`X^{},Y^0`$ could be produced in association with $`T`$ or $`Q`$, e.g., $`bTX^{}`$, these processes have two TeV-mass particles in the final state and will be limited by phase space. The production cross sections of the $`X`$ and $`Y`$ gauge bosons in Drell-Yan are very small. We thus consider other ways of producing these particles. If they are light enough, $`X`$ and $`Y`$ can be produced in the decays of the TeV-scale quark partners: $$TX^+b,\overline{Y}^0t,U_jX^+d_j,\overline{Y}^0u_j\mathrm{or}D_jX^{}u_j,Y^0d_j.$$ (52) For example, taking $`M_T=1`$ TeV, $`M_Y=0.9`$ TeV and $`\lambda _T=1`$, we find ($`Tt\overline{Y}^0`$ is kinematically forbidden for these masses), $$\mathrm{BR}(TbX^+)0.55\%.$$ (53) Similarly, $`X`$ and $`Y`$ can be produced through the decays of the heavy lepton partners, $`NX^+\mathrm{}^{},\overline{Y}^0\nu `$. The $`X`$ and $`Y`$ bosons can also be pair produced by electroweak interactions via the triple gauge couplings in Table 3; however, pair production of these TeV-scale particles will suffer from reduced phase space and off-shell $`s`$-channel propagators compared to Drell-Yan production of the $`Z^{}`$. If they are heavy enough, $`X`$ and $`Y`$ can decay to one SM fermion and one TeV-scale fermion partner, $`X^+T\overline{b},U_j\overline{d}_j,N_i\mathrm{}^+,Y^0t\overline{T},u_j\overline{U}_j,\nu _i\overline{N}_i(\mathrm{universal})`$ $`X^+T\overline{b},u_j\overline{D}_j,N_i\mathrm{}^+,Y^0t\overline{T},d_j\overline{D}_j,\nu _i\overline{N}_i(\mathrm{anomaly}\mathrm{free}).`$ (54) Neglecting the SM fermion mass, the partial widths for these decays are given by $$\mathrm{\Gamma }(VF\overline{f})=\frac{N_cg^2}{32\pi }\beta ^2\left[1\frac{\beta }{3}\right]M_V=4.2N_c\beta ^2\left[1\frac{\beta }{3}\right]\left(\frac{M_V}{\mathrm{TeV}}\right)\mathrm{GeV},$$ (55) where $`N_c=1`$ or 3 is the number of colors and $`\beta =(1M_F^2/M_V^2)`$. This decay mode and the production in Eq. (52) are mutually exclusive, depending on the relative masses of $`X,Y`$ and the TeV-scale fermion partners. If the decay to one SM fermion and one TeV-scale fermion partner is kinematically inaccessible, $`X`$ and $`Y`$ can decay to pairs of SM fermions through their mixings with the TeV-scale fermion partners, with partial widths proportional to $`\delta _t^2,\delta _\nu ^2v^2/f^2`$. The decays of $`X`$ are independent of the fermion embedding, $$X^{}b\overline{t},d_j\overline{u}_j,\mathrm{}^{}\overline{\nu },$$ (56) while the decays of $`Y`$ depend on the fermion embedding, since $`Y`$ can decay only to fermions that mix with a heavy partner: $$Y^0t\overline{t},u_j\overline{u}_j,\nu \overline{\nu }(\mathrm{universal}),Y^0t\overline{t},d_j\overline{d}_j,\nu \overline{\nu }(\mathrm{anomaly}\mathrm{free}).$$ (57) Unfortunately, there are no decays of $`Y`$ to charged dileptons because $`N_i`$ mix only with the neutrinos. The decays $`Y^0t\overline{t}`$, $`X^{}b\overline{t}`$ are controlled by $`\delta _t`$, while the decays to the first two quark generations and to the leptons are controlled by the smaller $`\delta _\nu `$. Thus, decays to third generation quarks will have a somewhat larger partial width. Neglecting final-state masses, the relevant partial widths are $`\mathrm{\Gamma }(X^{}b\overline{t})=\mathrm{\Gamma }(Y^0t\overline{t})`$ $`=`$ $`{\displaystyle \frac{3g^2}{48\pi }}\delta _t^2M_Y=0.51\lambda _T^2\left({\displaystyle \frac{\mathrm{TeV}}{M_T}}\right)^2\left({\displaystyle \frac{M_Y}{\mathrm{TeV}}}\right)\mathrm{GeV},`$ $`\mathrm{\Gamma }(X^{}jj)=\mathrm{\Gamma }(Y^0jj)`$ $`=`$ $`2{\displaystyle \frac{3g^2}{48\pi }}\delta _\nu ^2M_Y={\displaystyle \frac{0.11}{t_\beta ^2}}\left({\displaystyle \frac{\mathrm{TeV}}{M_Y}}\right)\mathrm{GeV},`$ $`\mathrm{\Gamma }(X^{}\mathrm{}\overline{\nu })=\mathrm{\Gamma }(Y^0\nu \overline{\nu })`$ $`=`$ $`3{\displaystyle \frac{g^2}{48\pi }}\delta _\nu ^2M_Y={\displaystyle \frac{0.054}{t_\beta ^2}}\left({\displaystyle \frac{\mathrm{TeV}}{M_Y}}\right)\mathrm{GeV},`$ (58) where $`jj`$ denote jets from quarks of the first two generations and the decays to leptons are summed over all three generations. Finally, $`Y`$ can decay to $`H\eta `$ via the coupling in the last row of Table 4, $$\mathrm{\Gamma }(Y^0H\eta )=\mathrm{\Gamma }(\overline{Y}^0H\eta )=\frac{g^2M_Y}{384\pi }=0.35\left(\frac{M_Y}{\mathrm{TeV}}\right)\mathrm{GeV}.$$ (59) ### 5.3 The singlet pseudoscalar $`\eta `$ The scalar sectors of little Higgs models are very model-dependent. For completeness, however, we briefly sketch here the decay modes of the singlet (pseudo-)scalar $`\eta `$ in the SU(3) simple group model. A more detailed analysis of the $`\eta `$ phenomenology can be found in Ref. . The singlet scalar $`\eta `$, which naturally gets a mass of a couple hundred GeV, can decay to pairs of SM fermions with couplings that depend on the SM fermion masses. These couplings receive contributions from the usual fermion Yukawa couplings, via the expansion of the nonlinear sigma model fields, and from the couplings of $`\eta `$ to a SM fermion and its TeV-scale partner combined with the $`F`$$`f`$ mixing. These couplings are all of order $`m_f/f`$, that is, suppressed by $`v/f`$ relative to the usual fermion Yukawa couplings. The $`\eta `$ can also decay into a Higgs boson and an off-shell $`Y`$, which then decays to a pair of SM fermions with couplings suppressed by the $`F`$$`f`$ mixing. We expect the decays of $`\eta `$ into pairs of fermions to dominate, with branching fractions proportional to the fermion masses up to order-one factors related to the contribution from the $`F`$$`f`$ mixing. The total width of $`\eta `$ will be suppressed by $`v^2/f^2`$ compared to that of a “bosophobic” Higgs of the same mass; however, this width will be too narrow to measure directly and too wide to give rise to displaced vertices, and thus can only be probed through production cross sections. ## 6 Conclusions The little Higgs models represent a new approach to electroweak symmetry breaking that will be accessible at future high-energy colliders. These models stabilize the hierarchy between a relatively low cutoff scale $`10`$ TeV and the electroweak scale by making the Higgs a pseudo-Goldstone boson of a spontaneously broken approximate global symmetry. Implementing such a global symmetry requires enlarging the gauge, fermion and scalar sectors of the SM. Little Higgs models therefore predict new gauge bosons, fermions and scalars at or below the TeV scale, which offer exciting possibilities for beyond-the-SM collider phenomenology at the LHC. However, many models of physics beyond the SM contain new gauge bosons, fermions, and/or scalars at or below the TeV scale. If such particles are discovered, one will want to know whether they implement the little Higgs mechanism by canceling the one-loop quadratic divergence in the Higgs mass due to the SM gauge bosons, top quark, and Higgs quartic coupling. We categorized the many little Higgs models into two classes based on the structure of the extended electroweak gauge group: * product group models, in which the SM SU(2)<sub>L</sub> gauge group arises from the diagonal breaking of two or more gauge groups, and * simple group models, in which the SM SU(2)<sub>L</sub> gauge group arises from the breaking of a single larger gauge group down to an SU(2) subgroup. As prototypes of each class, we studied the experimental signatures of the Littlest Higgs model and the SU(3) simple group model, respectively. The “smoking guns” for the little Higgs mechanism – the cancellation of the Higgs mass quadratic divergences between loops of SM particles and loops of the new particles – are quite straightforward and allow one to distinguish models that implement the little Higgs mechanism from other models that have a similar superficial phenomenology. In the top sector, the little Higgs mechanism appears as a sum rule involving the top quark Yukawa coupling, the $`TtH`$ or $`TbW`$ coupling $`\lambda _T`$, and the dimension-five $`TTHH`$ coupling $`\lambda _T^{}`$. In product group models, the simple structure of the top mass generation mechanism ensures that $`\lambda _T^{}`$ can be expressed in terms of $`\lambda _T`$, $`M_T`$ and the top Yukawa coupling. The little Higgs mechanism can then be checked by measuring $`\lambda _T`$ and $`M_T`$, computing the condensate $`f`$, and comparing with $`f`$ from the gauge sector. In simple group models, on the other hand, the top mass generation mechanism is slightly more complicated and involves two (or more) TeV-scale condensates. This introduces an extra free parameter into the top sector (which can be chosen as the ratio of the two condensates, $`f_2/f_1t_\beta `$), so that all three parameters $`\lambda _T`$, $`\lambda _T^{}`$, and $`M_T`$ must be measured in the top sector. We have not found a way to measure $`\lambda _T^{}`$ directly at the LHC. Instead, the required third parameter can be measured from the production rate of the TeV-scale quarks associated with the first two generations in the simple group models. These measurements of the extended top sector and the TeV-scale quark partners of the first two generations, if present, thus allow one to test the little Higgs mechanism in the top sector, distinguish the structure of the top quark mass generation mechanism, and extract the model parameters that control the fermion sector. We showed explicitly how these measurements allow one to distinguish the top sector of a little Higgs model from a fourth-generation top-prime and from a top see-saw model. In the gauge sector, the little Higgs mechanism appears as a sum rule involving the Higgs boson coupling to pairs of SM vector bosons and to pairs of the new TeV-scale vector bosons. The couplings involved in the sum rule can be directly measured via $`q\overline{q}V^{}V^{}H`$ associated production. Measurement of these couplings allows one to test which new particles are responsible for canceling each of the SM contributions to the Higgs mass-squared quadratic divergence. In product group models, the test of the little Higgs mechanism is particularly simple because of the collective breaking structure of the Higgs couplings to gauge bosons: it is enough to measure the $`Z_HZH`$ ($`W_HWH`$) couplings, which are accessible through $`Z_HZH`$ ($`W_HWH`$) decays. The simple group models have a different collective breaking structure in the gauge sector, however, so that a direct measurement of the $`V^{}V^{}H`$ couplings is necessary. Additional measurements in the gauge sector will shed light on the structure of the extended electroweak gauge group. We showed explicitly how measurements of the properties of a $`Z^{}`$ allow one to distinguish the $`Z^{}`$ states present in little Higgs models from the $`Z^{}`$s in the $`E_6`$ and left-right symmetric models and from a sequential $`Z^{}`$. The scalar sector is very model dependent. It depends on the global symmetry structure; therefore the classification of models into product group and simple group does not give a useful classification of the scalar sector phenomenology. ###### Acknowledgments. We thank B. Dobrescu, C. Hill, B. McElrath, J. Terning, D. Rainwater, M. Schmaltz, T. Tait, and W. Skiba for useful conversations. HEL and LTW thank the Aspen Center for Physics for hospitality while this work was initiated. This work was supported in part by the U.S. Department of Energy under grant DE-FG02-95ER40896 and in part by the Wisconsin Alumni Research Foundation. TH was also supported in part by the National Natural Science Foundation of China. LW was also supported in part by U.S. Department of Energy under grant DE-FG02-91ER40654. ## Appendix A Survey of little Higgs models ### A.1 Product group models The majority of little Higgs models are product group models. In addition to the Littlest Higgs, these include the theory space models (the Big Moose and the Minimal Moose ), the SU(6)/Sp(6) model of Ref. , and two extensions of the Littlest Higgs with built-in custodial SU(2) symmetry . There are also product group models with $`T`$-parity in the literature ; however, we do not address them here in any detail. In general, the phenomenology of models with $`T`$-parity is quite different from that discussed here; however, the top partner is typically $`T`$-parity even so that its phenomenology can be taken over directly from the Littlest Higgs case. We start with the theory space models. The Minimal Moose consists of two sites (where the gauge groups live) connected by four link fields (scalar fields transforming under the gauge groups at either end of the link). The electroweak gauge symmetry at one site is SU(2)$`\times `$U(1), while at the other it is SU(3) \[or alternatively, a second copy of SU(2)$`\times `$U(1); electroweak precision constraints favor this second possibility\]. The diagonal breaking of the gauge symmetry down to SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> leaves a set of SU(3) gauge bosons \[alternatively the broken SU(2)$`\times `$U(1) gauge bosons\] at the TeV scale. The top quark mass is generated by an interaction of the same form as Eq. (1), leaving a heavy charge 2/3 electroweak singlet quark at the TeV scale. The scalar spectrum consists of two Higgs doublets, a complex triplet and a complex singlet at the weak scale, with an additional Higgs doublet, triplet, and singlet at the TeV scale. The Big Moose is an extended version of this structure, with a longer chain of gauge groups connected by link fields that break down to the diagonal SU(2)$`\times `$U(1), leaving a larger number of broken gauge generators at the TeV scale. Many different theory space structures yield the little Higgs mechanism, with only mild topological constraints on the shape of the theory space . In particular, the theory space can be chosen such that the low-energy theory contains only two Higgs doublets, giving the extra light scalars of the Minimal Moose masses at the TeV scale . Theory space models always contain at least two light Higgs doublets. The SU(6)/Sp(6) model is similar to the Littlest Higgs, but starting with a global SU(6) symmetry broken down to Sp(6) at the TeV scale by an antisymmetric condensate. A subgroup \[SU(2)$`\times `$U(1)\]<sup>2</sup> of the global symmetry is gauged; the gauge symmetry is broken down to SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> by the condensate, leaving a set of SU(2)$`\times `$U(1) gauge bosons at the TeV scale. The top quark mass is generated in exact analogy to Eq. (1), leaving a heavy charge 2/3 electroweak singlet quark at the TeV scale. The scalar spectrum consists of two light Higgs doublets, plus a complex singlet at the TeV scale. The extensions of the Littlest Higgs with built-in custodial SU(2) symmetry were constructed in order to avoid some of the electroweak precision constraints on the Littlest Higgs model . The first such extension is a hybrid of the Littlest Higgs and the Minimal Moose with an SO(5)$`\times `$\[SU(2)$`\times `$U(1)\] gauge symmetry . It contains two light Higgs doublets, plus additional scalars at the TeV scale due to the enlarged global symmetry. It also contains extra TeV-scale gauge bosons from the enlarged gauge symmetry. The second such extension expands the global symmetry group to SO(9), spontaneously broken down to SO(5)$`\times `$SO(4) . This model contains only a single light Higgs doublet, with three scalar triplets and a singlet at the TeV scale. The gauge symmetry is \[SU(2)$`{}_{L}{}^{}\times `$SU(2)<sub>R</sub>\]$`\times `$\[SU(2)$`\times `$U(1)\], broken down to the SM electroweak gauge group by the symmetry breaking condensate. The model thus contains extra TeV-scale gauge bosons compared to the Littlest Higgs. The top sectors of both extensions are identical to that of the Littlest Higgs. The product group models all share two features. First, the models all contain a set of SU(2) gauge bosons at the TeV scale, obtained from the diagonal breaking of two gauge groups down to SU(2)<sub>L</sub>. Some models contain additional TeV-scale gauge bosons as well, from the breaking of more than two SU(2) gauge groups or from the breaking of gauge groups larger than SU(2). Second, the models all generate the top quark mass from a Lagrangian involving two terms, only one of which couples to the scalar sector of the model. This results in an extended top quark sector of the same form as in the Littlest Higgs model. These two features distinguish the product group models from the simple group models, which we consider next. ### A.2 Simple group models In addition to the SU(3) simple group model, there are two other simple group models in the literature to date: the SU(4) simple group model and the SU(9)/SU(8) model of Ref. . These two models depart from the SU(3) simple group model in different directions. The SU(4) simple group model is a straightforward extension of the SU(3) model to the electroweak gauge group SU(4)$`\times `$U(1)<sub>X</sub>. It was introduced because the simplest version of the SU(3) model generates a Higgs quartic coupling only at one-loop level through the Coleman-Weinberg potential, leading to a too-light Higgs boson . This problem can be fixed by adding an extra term to the scalar Lagrangian , which explicitly breaks a global U(1) symmetry in the model (and has the added benefit of giving mass to the $`\eta `$ pseudoscalar, which would otherwise be a Nambu-Goldstone boson). The SU(4) model, on the other hand, generates a Higgs quartic coupling at tree-level, so the Higgs mass is easily large enough. In the SU(4) simple group model the isospin doublets of the SM are all extended to quadruplets under SU(4). A total of four scalar quadruplets are needed to break SU(4)$`\times `$U(1)<sub>X</sub> down to SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub>, which leads to extra light scalars so that the low-energy theory contains two light Higgs doublets and two real singlets, plus three complex singlets which get masses of order $`f`$ TeV. The potential generated for the two Higgs doublets is not the most general possible, yielding interesting relations among the Higgs masses and couplings; in fact, the potential for the two Higgs doublets is of the same form as the one in the SU(6)/Sp(6) product group model. There are now four symmetry breaking vevs, $`f_{1,\mathrm{},4}`$. The fermion sector contains two heavy quark-partners and two heavy lepton-partners for each generation. Only one of the heavy quark-partners in each generation mixes with the corresponding SM quark. Like in the SU(3) model, the fermions can be embedded in a universal (but anomalous) way into SU(4) or in an anomaly-free way . Again, the anomaly-free embedding only works if the number of fermion generations is a multiple of three. The heavy gauge sector contains the broken generators of SU(4)$``$SU(2), namely two neutral gauge bosons $`Z^{}`$ and $`Z^{\prime \prime }`$ (which mix in general), two complex SU(2) doublets $`(Y^0,X^{})`$, $`(Y^0,X^{})`$, and a complex SU(2) singlet $`Y^{0\prime \prime }`$. The phenomenology of the first $`Z^{}`$ and the first doublet $`(Y^0,X^{})`$ are similar to those of the SU(3) model. The SU(9)/SU(8) model of Ref. contains exactly the same gauge group and fermion sector as the SU(3) simple group model. Thus the gauge and fermion sectors contain the same particle content and interactions as in the SU(3) simple group model. The only difference is the global symmetry structure, which leads to a different scalar sector. The global symmetry group is SU(9), broken down to SU(8) by a vacuum condensate with two independent vevs, $`f_{1,2}`$. The Higgs quartic coupling in this model is generated at tree level by Lagrangian terms that explicitly break the SU(9) global symmetry. The scalar sector contains two light Higgs doublets, plus two complex singlets that get masses of order $`f`$ TeV. As in the SU(4) model, the potential generated for the two Higgs doublets is far from the most general possible, yielding interesting relations among the Higgs masses and couplings. The simple group models share two features which distinguish them from the product group models. First, the models all contain an SU($`N`$)$`\times `$U(1) gauge symmetry that is broken down to SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub>, yielding the TeV-scale gauge bosons. The gauge couplings of the expanded SU($`N`$)$`\times `$U(1) symmetry are thus fixed in terms of the known SM gauge couplings. The gauge structure also forbids mixing between the SM $`W^\pm `$ bosons and the TeV-scale gauge bosons, in contrast to the product group models. Second, the top quark mass is generated from a Lagrangian involving two terms, which couple the top quark to two different nonlinear sigma model fields. This structure introduces an additional parameter into the top sector, which complicates the phenomenology and allows the heavy top-partner to be made lighter relative to the TeV-scale gauge bosons than in the product group models, thereby reducing the fine-tuning. ## Appendix B The SU(3) simple group model In this Appendix we collect some technical details of the SU(3) simple group model of Refs. and derive the interaction Lagrangian in the mass basis. The SU(3) simple group model is constructed by enlarging the SM SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> gauge group to SU(3)$`\times `$U(1)<sub>X</sub>. This requires enlarging the SU(2) doublets of the SM to SU(3) triplets and adding the additional SU(3) gauge bosons. The SU(3)$`\times `$U(1)<sub>X</sub> gauge symmetry is broken down to the SM electroweak gauge group by two complex scalar fields $`\mathrm{\Phi }_{1,2}`$, which are triplets under the SU(3) with aligned vevs $`f_{1,2}`$, both of order a TeV. We start with a scalar potential for $`\mathrm{\Phi }_{1,2}`$ which has a \[SU(3)$`\times `$U(1)\]<sup>2</sup> global symmetry. After $`\mathrm{\Phi }_{1,2}`$ acquire vevs, the global symmetry is spontaneously broken down to \[SU(2)$`\times `$U(1)\]<sup>2</sup>. At the same time, the global symmetry is broken explicitly down to its diagonal SU(3)$`\times `$U(1) subgroup by the gauge interactions. The scalar fields are parameterized as a nonlinear sigma model with $$\mathrm{\Theta }=\frac{1}{f}\left[\left(\begin{array}{cc}\begin{array}{cc}0& 0\\ 0& 0\end{array}& h\\ h^{}& 0\end{array}\right)+\frac{\eta }{\sqrt{2}}\left(\begin{array}{ccc}1& 0& \hfill 0\\ 0& 1& \hfill 0\\ 0& 0& \hfill 1\end{array}\right)\right],h=\left(\begin{array}{c}h^0\\ h^{}\end{array}\right),$$ (60) and $`\mathrm{\Phi }_1`$ $`=`$ $`e^{i\mathrm{\Theta }f_2/f_1}\left(\begin{array}{c}0\\ 0\\ f_1\end{array}\right)=fc_\beta \left[\left(\begin{array}{c}0\\ 0\\ 1\end{array}\right)+{\displaystyle \frac{it_\beta }{f}}\left(\begin{array}{c}h\\ \eta /\sqrt{2}\end{array}\right){\displaystyle \frac{t_\beta ^2}{2f^2}}\left(\begin{array}{c}\sqrt{2}\eta h\\ h^{}h+\eta ^2/2\end{array}\right)+\mathrm{}\right],`$ (71) $`\mathrm{\Phi }_2`$ $`=`$ $`e^{i\mathrm{\Theta }f_1/f_2}\left(\begin{array}{c}0\\ 0\\ f_2\end{array}\right)=fs_\beta \left[\left(\begin{array}{c}0\\ 0\\ 1\end{array}\right){\displaystyle \frac{i}{t_\beta f}}\left(\begin{array}{c}h\\ \eta /\sqrt{2}\end{array}\right){\displaystyle \frac{1}{2t_\beta ^2f^2}}\left(\begin{array}{c}\sqrt{2}\eta h\\ h^{}h+\eta ^2/2\end{array}\right)+\mathrm{}\right].`$ (82) We define $`f^2f_1^2+f_2^2`$ and $`t_\beta \mathrm{tan}\beta =f_2/f_1`$. Under the SU(2)<sub>L</sub> SM gauge group, $`h`$ transforms as a doublet and will be identified as the SM Higgs doublet with a vev $`v\sqrt{2}h^0=246`$ GeV, while $`\eta `$ is a real singlet which also remains light. We have chosen $`\eta `$ proportional to the unit matrix because this state remains unmixed with the unphysical (eaten) Goldstone bosons after EWSB.<sup>7</sup><sup>7</sup>7We thank Dave Rainwater for enlightening discussions on this point. We do not write down the Goldstone bosons that are eaten by the broken gauge generators. The SU(3) gauge bosons can be written in matrix form as $$A^aT^a=\frac{A^3}{2}\left(\begin{array}{ccc}1& & \\ & 1& \\ & & 0\end{array}\right)+\frac{A^8}{2\sqrt{3}}\left(\begin{array}{ccc}1& & \\ & 1& \\ & & 2\end{array}\right)+\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}& W^+& Y^0\\ W^{}& & X^{}\\ \overline{Y}^0& X^+& \end{array}\right).$$ (83) The $`\mathrm{\Phi }`$ vevs break the SU(3)$`\times `$U(1)<sub>X</sub> gauge symmetry down to the SM SU(2)$`{}_{L}{}^{}\times `$U(1)<sub>Y</sub> via the covariant derivative term $$_\mathrm{\Phi }=\left\|\left(_\mu +igA_\mu ^aT^a\frac{ig_x}{3}B_\mu ^x\right)\mathrm{\Phi }_i\right\|^2,$$ (84) where the SU(3) gauge coupling $`g`$ is equal to the SM SU(2)<sub>L</sub> gauge coupling and the U(1)<sub>X</sub> gauge coupling $`g_x`$ is fixed in terms of $`g`$ and the weak mixing angle $`t_W\mathrm{tan}\theta _W`$ by $$g_x=\frac{gt_W}{\sqrt{1t_W^2/3}}.$$ (85) The broken gauge generators get masses of order $`f`$ TeV and consist of a $`Z^{}`$ boson (a linear combination of $`A^8`$ and $`B^x`$) and a complex SU(2)<sub>L</sub> doublet $`(Y^0,X^{})`$. ### B.1 Gauge and Higgs sectors Before EWSB, the $`X`$ and $`Y`$ gauge bosons and a linear combination $`Z^{}`$ of the $`A^8`$ and $`B^x`$ gauge bosons get masses from the $`f`$ vevs. The linear combination $`Z^{}`$ that becomes massive is $$Z_0^{}=\frac{\sqrt{3}gA^8+g_xB^x}{\sqrt{3g^2+g_x^2}}=\frac{1}{\sqrt{3}}\left(\sqrt{3t_W^2}A^8+t_WB^x\right).$$ (86) We denote states and masses before EWSB with the subscript zero. The orthogonal combination of $`A^8`$ and $`B^x`$ becomes the hypercharge gauge boson $`B`$, $$B=\frac{g_xA^8+\sqrt{3}gB^x}{\sqrt{3g^2+g_x^2}}=\frac{1}{\sqrt{3}}\left(t_WA^8+\sqrt{3t_W^2}B^x\right).$$ (87) Hypercharge is given by $$Y=\frac{1}{\sqrt{3}}T^8+Q_x,T^8=\frac{1}{2\sqrt{3}}\mathrm{diag}(1,1,2),$$ (88) where $`Q_x=1/3`$ for the scalar fields $`\mathrm{\Phi }_i`$. We also have the relations $`A^3`$ $`=`$ $`c_WZ_0+s_WA,A^8=\sqrt{1t_W^2/3}Z_0^{}+{\displaystyle \frac{s_W^2}{\sqrt{3}c_W}}Z_0{\displaystyle \frac{s_W}{\sqrt{3}}}A`$ $`B^x`$ $`=`$ $`{\displaystyle \frac{t_W}{\sqrt{3}}}Z_0^{}s_W\sqrt{1t_W^2/3}Z_0+c_W\sqrt{1t_W^2/3}A,`$ (89) where $`A`$ is the photon. For use in precision corrections, we give the $`W`$ and $`Z`$ boson masses and their couplings to the Higgs at next-to-leading order in $`v^2/f^2`$ in Table 8. The $`WWH`$ and $`ZZH`$ couplings can be written in the form $$=2\frac{M_W^2}{v}y_WW^+W^{}H+\frac{M_Z^2}{v}y_ZZZH,$$ (90) with coefficients $`y_{W,Z}`$ given in Table 8. ### B.2 Fermion sector Because the model contains a gauged SU(3), SM fermions that are doublets under SU(2) must be expanded into triplets under the SU(3). In addition, new SU(3)-singlet fermions must be introduced to cancel the hypercharge anomalies and to marry and give mass to the new third components of the SU(3)-triplet fermions. The most straightforward way to construct a fermion sector for the SU(3) simple group model is to expand all the SU(2) doublets of the SM into SU(3) triplets, adding additional SU(3)-singlet right-handed fermions as needed, as was done in Ref. . We call this embedding “universal”, since the three generations have identical quantum numbers. The quarks and leptons of each generation are put into $`\mathrm{𝟑}`$ representations of SU(3): $`Q_m^T=(u,d,iU)_m,`$ $`iu_m^c,id_m^c,iU_m^c(\mathrm{universal})`$ $`L_m^T=(\nu ,e,iN)_m,`$ $`ie_m^c,iN_m^c,`$ (91) where $`m`$ is the generation index. We do not include a right-handed neutrino at this stage, leaving the neutrinos massless. Neutrino masses could be incorporated, e.g., through a see-saw mechanism in the UV completion of the little Higgs model or within the little Higgs theory itself ; however, this is beyond the scope of our current work. The $`Q_x`$ charges of the fermions are given in Table 9. It was pointed out by Kong that such a universal fermion sector leads to SU(3) and U(1)<sub>x</sub> gauge anomalies, although the SM SU(2) and U(1)<sub>Y</sub> gauge groups remain anomaly-free. These anomalies are not necessarily a problem because the little Higgs model is only an effective theory valid up to an energy scale $`\mathrm{\Lambda }4\pi f`$. Additional fermions can be added at the scale $`\mathrm{\Lambda }`$ to cancel the SU(3) and U(1)<sub>x</sub> gauge anomalies without affecting the phenomenology at the $`f`$ scale. Alternatively, one can construct a fermion sector that is anomaly-free already at the $`f`$ scale and yet contains no more degrees of freedom than the universal embedding, as proposed by Kong . This can be done by putting the first two generations of quarks in $`\overline{\mathrm{𝟑}}`$ representations of SU(3), while the third quark generation and all three lepton generations are in $`\mathrm{𝟑}`$s of SU(3). We call this embedding “anomaly-free”. It is fascinating to note that with this fermion content, the anomalies do not cancel within a single generation, as in the SM, but rather three generations (or a multiple thereof) are required to cancel the anomalies. The anomaly cancellation pattern of this fermion content has been previously pointed out in 3-3-1 models outside of the little Higgs context. The quarks of the third generation and three generations of leptons are put into $`\mathrm{𝟑}`$ representations of SU(3), exactly as in the universal embedding. The first two generations of quarks are put into $`\overline{\mathrm{𝟑}}`$ representations of SU(3): $`Q_1^T=(d,u,iD),`$ $`id^c,iu^c,iD^c(\mathrm{anomaly}\mathrm{free})`$ $`Q_2^T=(s,c,iS),`$ $`is^c,ic^c,iS^c,`$ (92) where the minus signs in front of $`u`$ and $`c`$ are there because the $`\overline{\mathrm{𝟐}}`$ of SU(2) is $`(d,u)`$ \[which is equivalent to the $`\mathrm{𝟐}`$, $`(u,d)`$\]. Notice that the heavy vector-like quarks of the first two generations have electric charge $`1/3`$, in contrast to the charge $`+2/3`$ heavy quark of the third generation. The $`Q_x`$ charges of the fermions are given in Table 9. #### B.2.1 Lepton masses and mixing The lepton sector is identical in both the universal and anomaly-free embeddings. The lepton masses are generated by the Lagrangian in Eq. (48), where we have chosen the flavor basis to correspond to the mass basis for the heavy neutrino partners $`N_m`$. The $`N_m`$ masses are then given by Eq. (49). The dimension-5 operator in Eq. (48) normalized by the cutoff scale $`\mathrm{\Lambda }`$ gives masses to the charged leptons via the $`3\times 3`$ Yukawa matrix $`\lambda _e^{mn}`$, which also generates a CKM-like mixing matrix $`V_{im}^{\mathrm{}}`$ between the charged lepton mass eigenstates $`e_i`$ and the heavy neutrino partners $`N_m`$. This mixing matrix appears in the $`X^{}\overline{e}_iN_m`$ couplings, $$\frac{g}{\sqrt{2}}V_{im}^{\mathrm{}}X_\mu ^{}\overline{e}_i\gamma ^\mu P_LN_m.$$ (93) These couplings can lead to lepton flavor violating processes, such as $`\mu e\gamma `$, via loops of $`N_m`$ and $`X^{}`$. As in the quark sector of the SM, this lepton flavor violation will be GIM-suppressed and will vanish in the limit that $`V_{im}^{\mathrm{}}`$ is diagonal, so that the $`N_m`$ mass eigenstates are aligned with the charged lepton mass eigenstates. The lepton flavor violation will also vanish in the limit that the $`N_m`$ are degenerate. The experimental limits on lepton flavor violation therefore put stringent constraints on the $`\lambda _{N_m}`$ couplings and/or on the structure of the $`\lambda _e^{mn}`$ matrix. After EWSB, the $`h`$ vev induces mixing between $`N_{m0}`$ and the corresponding neutrino $`\nu _{m0}`$ at order $`v/f`$, where as usual we use a subscript 0 to denote the SU(3) eigenstates and no subscript to denote the mass eigenstates after EWSB. Because of the structure of the $`N_m`$ mass term in Eq. (48), $`N_m`$ mixes only with the neutrino in the same SU(3) triplet, with a mixing angle $`\delta _\nu `$ given in Eq. (16) that is the same for all three generations. Note that $`t_\beta >1`$ suppresses $`\delta _\nu `$. The SU(3) eigenstates $`N_{m0}`$ and $`\nu _{i0}`$ are given in terms of the mass eigenstates $`N_m`$ and the SM neutrinos in the charged lepton mass basis ($`\nu _i=\nu _e,\nu _\mu ,\nu _\tau `$) by $$N_{m0}=N_m+\delta _\nu V_{mi}^{\mathrm{}}\nu _i,\nu _{i0}=\left(1\frac{1}{2}\delta _\nu ^2\right)\nu _i\delta _\nu V_{im}^{\mathrm{}}N_m,$$ (94) where we have kept the $`\delta _\nu ^2`$ term in the neutrino mixing because it will modify the well-measured couplings of neutrinos to the $`W`$ and $`Z`$ bosons at order $`v^2/f^2`$. In particular, the Fermi constant $`G_F`$ is measured in muon decay. The four-Fermi effective interaction Lagrangian is $$=2\sqrt{2}G_FJ^{+\mu }J_\mu ^{}=\frac{g^2}{2M_W^2}J^{+\mu }J_\mu ^{}\left(1\delta _\nu ^2\right).$$ (95) Plugging in $`M_W^2`$ (from Table 8) and $`\delta _\nu `$, we have, $$\frac{1}{G_F}=\sqrt{2}v^2\left\{1+\frac{v^2}{f^2}\left[\frac{1}{6}\left(\frac{s_\beta ^4}{c_\beta ^2}+\frac{c_\beta ^4}{s_\beta ^2}\right)+\frac{1}{2t_\beta ^2}\right]\right\}.$$ (96) The couplings of the scalars $`H`$ and $`\eta `$ to lepton pairs are given in Table 10. The couplings of charged leptons to $`H`$ get a multiplicative correction factor $`y_{\mathrm{}}`$ relative to the SM Yukawa couplings in terms of the lepton mass due to the nonlinear sigma model expansion. #### B.2.2 Lepton couplings to gauge bosons The fermion couplings to gauge bosons are given by the fermion kinetic term, $$=\overline{\psi }i𝒟_\mu \gamma ^\mu \psi ,𝒟=+igA^aT^a+ig_xQ_xB^x,$$ (97) with the $`Q_x`$ charges given in Table 9. The generators $`T^a`$ of the fundamental $`\mathrm{𝟑}`$ representation of SU(3) are given in Eq. (83). The couplings of the $`Z^{}`$ to lepton pairs were given in Table 4. The couplings of the heavy off-diagonal gauge bosons $`X^{}`$, $`Y^0`$ and $`\overline{Y}^0`$ to leptons were given in Table 2, neglecting flavor misalignment between the charged leptons and the $`N_m`$. Allowing for the possibility of flavor misalignment, we have $`_{X,Y}`$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}}}[iX_\mu ^{}\overline{e}_i\gamma ^\mu (V_{im}^{\mathrm{}}N_m+\delta _\nu \nu _i)+iY_\mu ^0\overline{\nu }_i\gamma ^\mu (V_{im}^{\mathrm{}}N_m+\delta _\nu \nu _i)+\mathrm{h}.\mathrm{c}.],`$ (98) where all fermion fields are left-handed and we have taken the neutrinos in the charged lepton mass basis, $`\nu _i=\nu _e,\nu _\mu ,\nu _\tau `$; $`N_m`$ are the heavy neutral leptons in their mass basis. The couplings of $`W^\pm `$ to lepton pairs, keeping terms of order $`v^2/f^2`$ in interactions involving only SM particles and terms of order $`v/f`$ in interactions involving one or more heavy particles, are $`_W`$ $`=`$ $`{\displaystyle \frac{gW_\mu ^+}{\sqrt{2}}}[(1{\displaystyle \frac{1}{2}}\delta _\nu ^2)\overline{\nu }_i\gamma ^\mu e_i\delta _\nu V_{mi}^{\mathrm{}}\overline{N}_m\gamma ^\mu e_i+\mathrm{h}.\mathrm{c}.].`$ (99) The couplings of the $`Z`$ boson to leptons, including the corrections from mixing between $`Z`$ and $`Z^{}`$ and mixing between the heavy neutral leptons and the SM neutrinos, are $`_Z`$ $`=Z_\mu {\displaystyle \frac{g}{c_W}}\{(J_3^\mu s_W^2J_Q^\mu ){\displaystyle \frac{1}{2}}\delta _\nu ^2\overline{\nu }_i\gamma ^\mu \nu _i{\displaystyle \frac{1}{2}}[\delta _\nu V_{im}^{\mathrm{}}\overline{N}_m\gamma ^\mu \nu _i+\mathrm{h}.\mathrm{c}.]`$ $`+`$ $`{\displaystyle \frac{\delta _Z}{\sqrt{34s_W^2}}}[({\displaystyle \frac{1}{2}}s_W^2)(\overline{\nu }_i\gamma ^\mu \nu _i+\overline{e}_i\gamma ^\mu e_i)+s_W^2\overline{e}_i^c\gamma ^\mu e_i^c+(1+s_W^2)\overline{N}_i\gamma ^\mu N_i]\},`$ where the leading-order coupling is given in terms of the standard fermion currents $$J_3^\mu =\overline{f}\gamma ^\mu T^3f,J_Q^\mu =\overline{f}\gamma ^\mu Q_ff\overline{f}^c\gamma ^\mu Q_{f^c}f^c.$$ (101) The couplings of the photon to fermions are given by the electromagnetic current as usual, $`_A=A_\mu eJ_Q^\mu `$. #### B.2.3 Quark masses and mixing: anomaly-free embedding The quark sector is more complicated than the lepton sector because of the anomaly-free embedding structure. The relevant Lagrangian terms for the third generation and for the first two generations are $`_3`$ $`=`$ $`\lambda _1^tiu_1^c\mathrm{\Phi }_1^{}Q_3+\lambda _2^tiu_2^c\mathrm{\Phi }_2^{}Q_3+{\displaystyle \frac{\lambda _b^m}{\mathrm{\Lambda }}}id_m^cϵ_{ijk}\mathrm{\Phi }_1^i\mathrm{\Phi }_2^jQ_3^k+\mathrm{h}.\mathrm{c}.`$ $`_{1,2}`$ $`=`$ $`\lambda _1^{dn}id_1^{nc}Q_n^T\mathrm{\Phi }_1+\lambda _2^{dn}id_2^{nc}Q_n^T\mathrm{\Phi }_2+{\displaystyle \frac{\lambda _u^{mn}}{\mathrm{\Lambda }}}iu_m^cϵ_{ijk}\mathrm{\Phi }_1^i\mathrm{\Phi }_2^jQ_n^k+\mathrm{h}.\mathrm{c}.,`$ (102) where $`n=1,2`$; $`i,j,k=1,2,3`$ are SU(3) indexes; $`u_1^c`$ and $`u_2^c`$ are linear combinations of $`t^c`$ and $`T^c`$ \[see Eqs. (103) and (105) below\]; $`b_m^c`$ runs over all the down-type conjugate quarks ($`d^c,s^c,b^c,D^c,S^c`$); $`d_1^{nc}`$ and $`d_2^{nc}`$ are linear combinations of $`d^c`$ and $`D^c`$ for $`n=1`$ and of $`s^c`$ and $`S^c`$ for $`n=2`$ \[see Eqs. (106) and (108) below\]; and $`u_m^c`$ runs over all the up-type conjugate quarks ($`u^c,c^c,t^c,T^c`$). The $`f`$ vevs generate mass terms for three heavy quarks. The state $$T^c=\frac{\lambda _1^tc_\beta u_1^c+\lambda _2^ts_\beta u_2^c}{\sqrt{\lambda _1^{t2}c_\beta ^2+\lambda _2^{t2}s_\beta ^2}}$$ (103) marries $`T`$, giving it a mass of $$M_T=f\sqrt{\lambda _1^{t2}c_\beta ^2+\lambda _2^{t2}s_\beta ^2}$$ (104) and leaving the orthogonal combination of $`u_1^c`$ and $`u_2^c`$ massless: $$t^c=\frac{\lambda _2^ts_\beta u_1^c+\lambda _1^tc_\beta u_2^c}{\sqrt{\lambda _1^{t2}c_\beta ^2+\lambda _2^{t2}s_\beta ^2}}.$$ (105) The states (here we denote $`\lambda _{1,2}^{dn}`$ by $`\lambda _{1,2}^d`$ for $`n=1`$ and by $`\lambda _{1,2}^s`$ for $`n=2`$) $$D^c=\frac{\lambda _1^dc_\beta d_1^{1c}+\lambda _2^ds_\beta d_2^{1c}}{\sqrt{\lambda _1^{d2}c_\beta ^2+\lambda _2^{d2}s_\beta ^2}},S^c=\frac{\lambda _1^sc_\beta d_1^{2c}+\lambda _2^ss_\beta d_2^{2c}}{\sqrt{\lambda _1^{s2}c_\beta ^2+\lambda _2^{s2}s_\beta ^2}}$$ (106) marry $`D`$ and $`S`$, respectively, giving them masses of $$M_D=f\sqrt{\lambda _1^{d2}c_\beta ^2+\lambda _2^{d2}s_\beta ^2},M_S=f\sqrt{\lambda _1^{s2}c_\beta ^2+\lambda _2^{s2}s_\beta ^2},$$ (107) and leaving the orthogonal combinations massless: $$d^c=\frac{\lambda _2^ds_\beta d_1^{1c}+\lambda _1^dc_\beta d_2^{1c}}{\sqrt{\lambda _1^{d2}c_\beta ^2+\lambda _2^{d2}s_\beta ^2}},s^c=\frac{\lambda _2^ss_\beta d_1^{2c}+\lambda _1^sc_\beta d_2^{2c}}{\sqrt{\lambda _1^{s2}c_\beta ^2+\lambda _2^{s2}s_\beta ^2}}.$$ (108) After EWSB, the quark mass terms are $`_{\mathrm{up}\mathrm{mass}}`$ $`=`$ $`M_TT^cT+{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{s_\beta c_\beta (\lambda _1^{t2}\lambda _2^{t2})}{\sqrt{\lambda _1^{t2}c_\beta ^2+\lambda _2^{t2}s_\beta ^2}}}T^ct{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{\lambda _1^t\lambda _2^t}{\sqrt{\lambda _1^{t2}c_\beta ^2+\lambda _2^{t2}s_\beta ^2}}}t^ct`$ (109) $`+{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{f}{\mathrm{\Lambda }}}\lambda _u^{mn}u_m^cu_n+\mathrm{h}.\mathrm{c}.`$ $`_{\mathrm{down}\mathrm{mass}}`$ $`=`$ $`M_DD^cD{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{s_\beta c_\beta (\lambda _1^{d2}\lambda _2^{d2})}{\sqrt{\lambda _1^{d2}c_\beta ^2+\lambda _2^{d2}s_\beta ^2}}}D^cd+{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{\lambda _1^d\lambda _2^d}{\sqrt{\lambda _1^{d2}c_\beta ^2+\lambda _2^{d2}s_\beta ^2}}}d^cd`$ (110) $`M_SS^cS{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{s_\beta c_\beta (\lambda _1^{s2}\lambda _2^{s2})}{\sqrt{\lambda _1^{s2}c_\beta ^2+\lambda _2^{s2}s_\beta ^2}}}S^cs+{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{\lambda _1^s\lambda _2^s}{\sqrt{\lambda _1^{s2}c_\beta ^2+\lambda _2^{s2}s_\beta ^2}}}s^cs`$ $`+{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{f}{\mathrm{\Lambda }}}\lambda _b^md_m^cb+\mathrm{h}.\mathrm{c}.`$ where $`u_n=u,c`$; $`u_m^c=u^c,c^c,t^c,T^c`$; and $`d_m^c=d^c,s^c,b^c,D^c,S^c`$. The couplings $`\lambda _u^{mn}`$ and $`\lambda _b^m`$ cause a misalignment between the mass eigenstates in the up and down sectors, leading to the CKM matrix. They also cause an analogous misalignment between the SM quark mass eigenstates and the heavy quarks $`D`$, $`S`$, and $`T`$, leading to an analogous matrix. We choose the “flavor basis” to be the mass basis for $`D,S,T`$. Two unitary matrices are needed to rotate the left-handed up- and down-type quarks from the flavor basis (primed fields) into the mass basis (unprimed fields): $$V^u\left(\begin{array}{c}u^{}\\ c^{}\\ t^{}\end{array}\right)=\left(\begin{array}{c}u\\ c\\ t\end{array}\right),V^d\left(\begin{array}{c}d^{}\\ s^{}\\ b^{}\end{array}\right)=\left(\begin{array}{c}d\\ s\\ b\end{array}\right).$$ (111) The CKM matrix is then given by $$V^{\mathrm{CKM}}=V^uV^d.$$ (112) These matrices appear in the quark gauge couplings; see Sec. B.2.4 for details. Note that, in contrast to the SM, there are two physically meaningful mixing matrices. Electroweak symmetry breaking also induces mixing between the heavy left-handed quarks $`D,S,T`$ and the SM quarks. In the up-quark sector, the terms in Eq. (109) involving $`T^c`$ lead to mixing between $`T`$ and $`u,c,t`$ that violates the SU(3) symmetry. As usual we use the subscript $`0`$ to denote SU(3) states; fields with no subscript denote the mass eigenstates after the mixing induced by EWSB. We can rewrite the SU(3) state $`T_0`$ in terms of the mass eigenstate $`T`$ and the SM fermions in the interaction basis (primed fields) as $$T_0=T+\delta _{u_i}u_i^{},$$ (113) with $`i=1,2,3`$, where $$\delta _u=\frac{v}{\sqrt{2}\mathrm{\Lambda }}\frac{\lambda _u^{T^cu}}{\sqrt{\lambda _1^{t2}c_\beta ^2+\lambda _2^{t2}s_\beta ^2}},\delta _c=\frac{v}{\sqrt{2}\mathrm{\Lambda }}\frac{\lambda _u^{T^cc}}{\sqrt{\lambda _1^{t2}c_\beta ^2+\lambda _2^{t2}s_\beta ^2}},\delta _t=\frac{v}{\sqrt{2}f}\frac{s_\beta c_\beta (\lambda _1^{t2}\lambda _2^{t2})}{(\lambda _1^{t2}c_\beta ^2+\lambda _2^{t2}s_\beta ^2)}.$$ (114) One can choose the couplings $`\lambda _u^{T^cu}`$ and $`\lambda _u^{T^cc}`$ to be small in order to suppress the mixing effects in the first and second generations. In the mass basis (unprimed fields) this becomes $$T_0=T+\mathrm{\Delta }_{u_i}u_i,\mathrm{\Delta }_{u_i}=V_{ij}^u\delta _{u_j}V_{i3}^u\delta _t,$$ (115) where in the last approximate equality we neglect $`\lambda _u^{T^cu}`$ and $`\lambda _u^{T^cc}`$. After mixing, the up quarks in the mass basis become $$u_{i0}=\left(1\frac{1}{2}\|\mathrm{\Delta }_{u_i}\|^2\right)u_i\mathrm{\Delta }_{u_i}T,$$ (116) where we have kept the $`\|\mathrm{\Delta }_{u_i}\|^2`$ term (which is of order $`v^2/f^2`$) because it will modify the well-measured couplings of quarks to the $`W`$ boson. Similarly, in the down-quark sector, the terms in Eq. (110) involving $`D^c`$ ($`S^c`$) lead to mixing between $`D`$ and $`d,b`$ ($`S`$ and $`s,b`$). As in the up sector, we can rewrite the SU(3) states $`D_0`$ and $`S_0`$ in terms of the mass eigenstates $`D`$ and $`S`$ and the SM fermions in the interaction basis (primed fields) as $$D_0=D+\delta _{Dd_i}d_i^{},S_0=S+\delta _{Sd_i}d_i^{},$$ (117) with $`i=1,2,3`$, where $`\delta _{Dd}`$ $`=`$ $`{\displaystyle \frac{v}{\sqrt{2}f}}{\displaystyle \frac{s_\beta c_\beta (\lambda _1^{d2}\lambda _2^{d2})}{(\lambda _1^{d2}c_\beta ^2+\lambda _2^{d2}s_\beta ^2)}},\delta _{Ds}=0,\delta _{Db}={\displaystyle \frac{v}{\sqrt{2}\mathrm{\Lambda }}}{\displaystyle \frac{\lambda _b^{D^c}}{\sqrt{\lambda _1^{d2}c_\beta ^2+\lambda _2^{d2}s_\beta ^2}}},`$ $`\delta _{Sd}`$ $`=`$ $`0,\delta _{Ss}={\displaystyle \frac{v}{\sqrt{2}f}}{\displaystyle \frac{s_\beta c_\beta (\lambda _1^{s2}\lambda _2^{s2})}{(\lambda _1^{s2}c_\beta ^2+\lambda _2^{s2}s_\beta ^2)}},\delta _{Sb}={\displaystyle \frac{v}{\sqrt{2}\mathrm{\Lambda }}}{\displaystyle \frac{\lambda _b^{S^c}}{\sqrt{\lambda _1^{s2}c_\beta ^2+\lambda _2^{s2}s_\beta ^2}}}.`$ (118) The zero mixings, $`\delta _{Ds}=\delta _{Sd}=0`$, are a consequence of the collective breaking mass generation for $`d`$ and $`s`$ in the $`D,S`$ mass basis. One can choose $`\lambda _b^{D^c}`$ and $`\lambda _b^{S^c}`$ to be small in order to suppress the mixing effects in the $`b`$ quark sector. From Eq. (110), the small mass of the $`d`$ ($`s`$) quark requires one of the couplings $`\lambda _{1,2}^d`$ ($`\lambda _{1,2}^s`$) to be very small. We choose the small coupling to be $`\lambda _1^d`$ ($`\lambda _1^s`$) so that the mixing effects in the down-quark sector are suppressed in the same $`t_\beta >1`$ limit as the mixing effects in the neutrino sector. We then have, $$\delta _{Dd}\delta _{Ss}\frac{v}{\sqrt{2}t_\beta f}=\delta _\nu .$$ (119) In the mass basis (unprimed fields), the $`D`$ and $`S`$ states become $`D_0=D+\mathrm{\Delta }_{Dd_i}d_i,`$ $`\mathrm{\Delta }_{Dd_i}=V_{ij}^d\delta _{Dd_j}V_{i1}^d\delta _\nu ,`$ (120) $`S_0=S+\mathrm{\Delta }_{Sd_i}d_i,`$ $`\mathrm{\Delta }_{Sd_i}=V_{ij}^d\delta _{Sd_j}V_{i2}^d\delta _\nu ,`$ where in the last approximate equalities we neglect $`\lambda _b^{D^c}`$ and $`\lambda _b^{S^c}`$. After mixing, the down quarks in the mass basis become $$d_{i0}=\left(1\frac{1}{2}\|\mathrm{\Delta }_{Dd_i}\|^2\frac{1}{2}\|\mathrm{\Delta }_{Sd_i}\|^2\right)d_i\mathrm{\Delta }_{Dd_i}D\mathrm{\Delta }_{Sd_i}S,$$ (121) where we again have kept the $`\|\mathrm{\Delta }_{Dd_i}\|^2`$ and $`\|\mathrm{\Delta }_{Sd_i}\|^2`$ terms, which are of order $`v^2/f^2`$. We now write the couplings of the scalars, $`H`$ and $`\eta `$, to quark pairs, taking into account corrections from the expansion of the nonlinear sigma model and the mixing between the SM quarks and the heavy quarks. The different treatment of the third quark generation in the anomaly-free fermion embedding \[Eq. (102)\] leads to flavor-changing couplings of quarks to $`H`$ (at order $`v^2/f^2`$) and to $`\eta `$ (at order $`v/f`$). The full parameter dependence of the flavor changing couplings depends on the exact form of the up and down quark mass matrices, which determine the quark mixing in the left- and right-handed sectors. A detailed exploration of the quark mass matrices is beyond the scope of this work. Instead, we write down the scalar couplings ignoring the mixing of the right-handed top quark $`t^c`$ with the first two generations. We begin with the couplings of $`T`$ quark pairs. $`T`$ couples to $`\eta `$ with a coupling of order one and to $`H`$ with a coupling of order $`v/f`$: $`_{T^cT}`$ $``$ $`(HT^cT){\displaystyle \frac{v}{f}}\left[(\lambda _1^{t2}s_\beta ^2+\lambda _2^{t2}c_\beta ^2){\displaystyle \frac{f}{2M_T}}s_\beta ^2c_\beta ^2(\lambda _1^{t2}\lambda _2^{t2})^2{\displaystyle \frac{f^3}{2M_T^3}}\right]`$ (122) $`+(i\eta T^cT)s_\beta c_\beta (\lambda _1^{t2}\lambda _2^{t2}){\displaystyle \frac{f}{\sqrt{2}M_T}}+\mathrm{h}.\mathrm{c}.,`$ where we have neglected terms involving $`\lambda _u^{T^cu}`$ and $`\lambda _u^{T^cc}`$. Similarly, $`D`$ and $`S`$ quark pairs couple to $`\eta `$ with a coupling of order one: $$_{D_m^cD_m}\frac{c_\beta }{\sqrt{2}}\lambda _2^d(i\eta D^cD)+\frac{c_\beta }{\sqrt{2}}\lambda _2^s(i\eta S^cS)+\mathrm{h}.\mathrm{c}.,$$ (123) where we have neglected terms involving $`\lambda _b^{D^c}`$ and $`\lambda _b^{S^c}`$ and taken $`\lambda _1^{d,s}\lambda _2^{d,s}`$ \[if the top quark mass were neglected, Eq. (122) would also reduce to this simple form\]. One would naively expect an $`HD_m^cD_m`$ coupling at order $`v/f`$ coming from replacing one Higgs field by its vev in the nonlinear sigma model expansion term $`HHD_m^cD_m`$; however, this term is exactly canceled by the contribution from $`HD_m^cd_m`$ after $`dD`$ mixing if the down and strange quark masses are neglected. The leading-order couplings of scalars to one $`T`$ quark and one SM up-type quark are $$(HT^cu_i)\left[s_\beta c_\beta (\lambda _1^{t2}\lambda _2^{t2})\frac{f}{\sqrt{2}M_T}V_{i3}^u\right](i\eta t^cT)\left[\frac{\lambda _1^t\lambda _2^tf}{\sqrt{2}M_T}\right]+\mathrm{h}.\mathrm{c}.,$$ (124) where we again neglect terms involving $`\lambda _u^{T^cu}`$ and $`\lambda _u^{T^cc}`$ and in the last term ignore the mixing of the right-handed top quark $`t^c`$ with the first two generations. The last term can be written in terms of SM quark masses and mixing angles via the relation (again ignoring right-handed quark mixing) $$\frac{\lambda _1^t\lambda _2^tf}{\sqrt{2}M_T}=\underset{j}{}\frac{m_{u_j}}{v}V_{j3}^u\frac{m_t}{v}V_{33}^u,$$ (125) where we have used $`m_u,m_cm_t`$. The couplings in Eq. (124) will lead to the decays $`TtH`$ and $`Tt\eta `$. The couplings of scalars to $`D,S`$ and one SM down-type quark are $$\frac{c_\beta }{\sqrt{2}}\lambda _2^dV_{i1}^d(HD^cd_i)+\frac{c_\beta }{\sqrt{2}}\lambda _2^sV_{i2}^d(HS^cd_i)+\mathrm{h}.\mathrm{c}.,$$ (126) where we have neglected terms involving $`\lambda _b^{D^c}`$ and $`\lambda _b^{S^c}`$ and ignored couplings of $`\eta `$ proportional to the down or strange quark masses. These couplings will lead to the decays $`D,Sd_iH`$. The couplings of scalars to a pair of SM up-type quarks (again ignoring right-handed quark mixing) are $``$ $`=`$ $`(Hu_i^cu_j)\{\delta _{ij}{\displaystyle \frac{m_{u_i}}{v}}[1{\displaystyle \frac{v^2}{6f^2}}(3+{\displaystyle \frac{s_\beta ^4}{c_\beta ^2}}+{\displaystyle \frac{c_\beta ^4}{s_\beta ^2}})]+\delta _{i3}{\displaystyle \frac{m_t}{2\sqrt{2}f}}V_{33}^u\mathrm{\Delta }_{u_j}\left({\displaystyle \frac{c_\beta ^2s_\beta ^2}{s_\beta c_\beta }}\right)`$ (127) $`\delta _{i3}V_{j3}^u{\displaystyle \frac{vm_t}{2f^2}}V_{33}^u+\delta _{u_i^c}\mathrm{\Delta }_{u_j}{\displaystyle \frac{M_T}{v}}\}`$ $`+(i\eta u_i^cu_j)\left[\delta _{ij}{\displaystyle \frac{m_{u_i}}{\sqrt{2}f}}\left({\displaystyle \frac{s_\beta ^2c_\beta ^2}{s_\beta c_\beta }}\right)+\delta _{i3}\mathrm{\Delta }_{u_j}{\displaystyle \frac{m_t}{v}}V_{33}^u\right]+\mathrm{h}.\mathrm{c}.`$ Note the flavor-changing couplings involving $`t^c`$ from terms containing a $`\delta _{i3}`$. Here we have introduced the notation $`\delta _{u_i^c}`$ for the mixings between $`u_i^c`$ and $`T^c`$, which occur at order $`v^2/f^2`$. They are given explicitly by $`\delta _{u^c}`$ $`=`$ $`{\displaystyle \frac{v}{M_T}}{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\left(\delta _u\lambda _u^{u^cu^{}}+\delta _c\lambda _u^{u^cc^{}}\right),\delta _{c^c}={\displaystyle \frac{v}{M_T}}{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\left(\delta _u\lambda _u^{c^cu^{}}+\delta _c\lambda _u^{c^cc^{}}\right),`$ $`\delta _{t^c}`$ $`=`$ $`{\displaystyle \frac{v}{M_T}}\left\{{\displaystyle \frac{m_t}{v}}V_{33}^u\left[\mathrm{\Delta }_t+{\displaystyle \frac{v}{2\sqrt{2}f}}\left({\displaystyle \frac{c_\beta ^2s_\beta ^2}{s_\beta c_\beta }}\right)\right]{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\left(\delta _u\lambda _u^{t^cu^{}}+\delta _c\lambda _u^{t^cc^{}}\right)\right\},`$ (128) where $`T^c=T_0^c\delta _{u_i^c}u_{i0}^c`$. The couplings of scalars to a pair of SM down-type quarks (again ignoring right-handed quark mixing) are $``$ $`=`$ $`(Hd_i^cd_j)\{\delta _{ij}{\displaystyle \frac{m_{d_i}}{v}}[1{\displaystyle \frac{v^2}{6f^2}}(3+{\displaystyle \frac{s_\beta ^4}{c_\beta ^2}}+{\displaystyle \frac{c_\beta ^4}{s_\beta ^2}})]+{\displaystyle \frac{v^2}{2f^2}}[\delta _{i1}{\displaystyle \frac{m_d}{v}}V_{11}^dV_{j1}^d\delta _{i2}{\displaystyle \frac{m_s}{v}}V_{22}^dV_{j2}^d]`$ (129) $`+{\displaystyle \frac{v}{2\sqrt{2}f}}\left({\displaystyle \frac{c_\beta ^2s_\beta ^2}{s_\beta c_\beta }}\right)\left[\delta _{i1}\mathrm{\Delta }_{Dd_j}{\displaystyle \frac{m_d}{v}}V_{11}^d\delta _{i2}\mathrm{\Delta }_{Sd_j}{\displaystyle \frac{m_s}{v}}V_{22}^d\right]`$ $`+\delta _{Dd_i^c}\mathrm{\Delta }_{Dd_j}{\displaystyle \frac{M_D}{v}}+\delta _{Sd_i^c}\mathrm{\Delta }_{Sd_j}{\displaystyle \frac{M_S}{v}}\}`$ $`+(i\eta d_i^cd_j)\left[\delta _{ij}{\displaystyle \frac{m_{d_i}}{\sqrt{2}f}}\left({\displaystyle \frac{s_\beta ^2c_\beta ^2}{s_\beta c_\beta }}\right)+\delta _{i1}\mathrm{\Delta }_{Dd_j}{\displaystyle \frac{m_d}{v}}V_{11}^d+\delta _{i2}\mathrm{\Delta }_{Sd_j}{\displaystyle \frac{m_s}{v}}V_{22}^d\right]+\mathrm{h}.\mathrm{c}.,`$ where we have used (neglecting right-handed quark mixing) $$\frac{\lambda _1^d\lambda _2^df}{\sqrt{2}M_D}=\frac{m_d}{v}V_{11}^d,\frac{\lambda _1^s\lambda _2^sf}{\sqrt{2}M_S}=\frac{m_s}{v}V_{22}^d.$$ (130) Note the flavor-changing couplings involving $`d^c`$ ($`s^c`$) from terms containing a $`\delta _{i1}`$ ($`\delta _{i2}`$). We also introduce the notation $`\delta _{Dd_i^c}`$, $`\delta _{Sd_i^c}`$ for the mixings between $`d_i^c`$ and $`D^c`$, $`S^c`$, respectively, which occur at order $`v^2/f^2`$. They are given explicitly by $`\delta _{Dd^c}`$ $`=`$ $`{\displaystyle \frac{v}{M_D}}\left\{{\displaystyle \frac{m_d}{v}}V_{11}^d\left[\delta _{Dd}{\displaystyle \frac{v}{2\sqrt{2}f}}\left({\displaystyle \frac{c_\beta ^2s_\beta ^2}{s_\beta c_\beta }}\right)\right]+{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\delta _{Db}\lambda _b^{d^c}\right\},`$ $`\delta _{Ds^c}`$ $`=`$ $`{\displaystyle \frac{v}{M_D}}{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\delta _{Db}\lambda _b^{s^c},\delta _{Db^c}={\displaystyle \frac{v}{M_D}}{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\delta _{Db}\lambda _b^{b^c},`$ $`\delta _{Ss^c}`$ $`=`$ $`{\displaystyle \frac{v}{M_S}}\left\{{\displaystyle \frac{m_s}{v}}V_{22}^d\left[\delta _{Ss}{\displaystyle \frac{v}{2\sqrt{2}f}}\left({\displaystyle \frac{c_\beta ^2s_\beta ^2}{s_\beta c_\beta }}\right)\right]+{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\delta _{Sb}\lambda _b^{s^c}\right\},`$ $`\delta _{Sd^c}`$ $`=`$ $`{\displaystyle \frac{v}{M_S}}{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\delta _{Sb}\lambda _b^{d^c},\delta _{Sb^c}={\displaystyle \frac{v}{M_S}}{\displaystyle \frac{f}{\sqrt{2}\mathrm{\Lambda }}}\delta _{Sb}\lambda _b^{b^c},`$ (131) where $`D^c=D_0^c\delta _{Dd_i^c}d_{i0}^c`$ and $`S^c=S_0^c\delta _{Sd_i^c}d_{i0}^c`$. #### B.2.4 Quark couplings to gauge bosons: anomaly-free embedding The couplings of the heavy off-diagonal gauge bosons $`X^{}`$, $`Y^0`$ and $`\overline{Y}^0`$ to quarks in the anomaly-free embedding were given in Table 4, neglecting flavor misalignment and CKM mixing. Allowing for the flavor misalignment, we have<sup>8</sup><sup>8</sup>8The SU(3) generators for the quarks of the first two generations, in the antifundamental $`\overline{\mathrm{𝟑}}`$ representation, are given by $`T^a`$. $`_{X,Y}`$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}}}\{iX_\mu ^{}\overline{d}_i\gamma ^\mu [V_{i3}^dT+(\mathrm{\Delta }_{u_j}V_{i3}^d+\mathrm{\Delta }_{Dd_i}^{}V_{j1}^u+\mathrm{\Delta }_{Sd_i}^{}V_{j2}^u)u_j]`$ (132) $`+iX_\mu ^+\overline{u}_i\gamma ^\mu V_{ij}^uD_j+iY_\mu ^0\overline{u}_i\gamma ^\mu \left(V_{i3}^uT+\mathrm{\Delta }_{u_k}V_{i3}^uu_k\right)`$ $`+i\overline{Y}_\mu ^0\overline{d}_i\gamma ^\mu [V_{ij}^dD_j+(\mathrm{\Delta }_{Dd_k}V_{i1}^d+\mathrm{\Delta }_{Sd_k}V_{i2}^d)d_k]+\mathrm{h}.\mathrm{c}.\}.`$ The couplings of $`W^\pm `$ to quark pairs, keeping terms of order $`v^2/f^2`$ in interactions involving only SM particles and terms of order $`v/f`$ in interactions involving one or more heavy particles, are $`_W`$ $`=`$ $`{\displaystyle \frac{gW_\mu ^+}{\sqrt{2}}}[(1{\displaystyle \frac{1}{2}}\|\mathrm{\Delta }_{u_i}\|^2{\displaystyle \frac{1}{2}}\|\mathrm{\Delta }_{Dd_j}\|^2{\displaystyle \frac{1}{2}}\|\mathrm{\Delta }_{Sd_j}\|^2)V_{ij}^{\mathrm{CKM}}\overline{u}_i\gamma ^\mu d_j`$ (133) $`V_{ij}^{\mathrm{CKM}}\mathrm{\Delta }_{u_i}^{}\overline{T}\gamma ^\mu d_jV_{ij}^{\mathrm{CKM}}\mathrm{\Delta }_{Dd_j}\overline{u}_i\gamma ^\mu DV_{ij}^{\mathrm{CKM}}\mathrm{\Delta }_{Sd_j}\overline{u}_i\gamma ^\mu S+\mathrm{h}.\mathrm{c}.].`$ The couplings of the $`Z^{}`$ boson to quarks were also given in Table 4, neglecting flavor misalignment and CKM mixing. Allowing for the flavor misalignment, we find flavor-changing couplings for the left-handed quarks involving $`V_{i3}^uV_{3j}^u`$ in the up sector and $`V_{i3}^dV_{3j}^d`$ in the down sector: $`_Z^{}`$ $``$ $`{\displaystyle \frac{g}{c_W}}{\displaystyle \frac{Z_\mu ^{}}{\sqrt{34s_W^2}}}[({\displaystyle \frac{1}{2}}+{\displaystyle \frac{2}{3}}s_W^2)(\overline{u}_i\gamma ^\mu u_i+\overline{d}_i\gamma ^\mu d_i)`$ (134) $`+(1s_W^2)(V_{i3}^uV_{3j}^u\overline{u}_i\gamma ^\mu u_j+V_{i3}^dV_{3j}^d\overline{d}_i\gamma ^\mu d_j)].`$ The couplings of the $`Z`$ boson to quarks, including the corrections from mixing between $`Z`$ and $`Z^{}`$ and mixing between the TeV-scale quarks and their SM partners, are $`_Z`$ $`=`$ $`Z_\mu {\displaystyle \frac{g}{c_W}}\{(J_3^\mu s_W^2J_Q^\mu )+{\displaystyle \frac{1}{2}}[\|\mathrm{\Delta }_{u_i}\|^2\overline{u}_i\gamma ^\mu u_i+(\|\mathrm{\Delta }_{Dd_i}\|^2+\|\mathrm{\Delta }_{Sd_i}\|^2)\overline{d}_i\gamma ^\mu d_i]`$ (135) $`+{\displaystyle \frac{\delta _Z}{\sqrt{34s_W^2}}}[({\displaystyle \frac{1}{2}}+{\displaystyle \frac{2}{3}}s_W^2)(\overline{u}_i\gamma ^\mu u_i+\overline{d}_i\gamma ^\mu d_i){\displaystyle \frac{2}{3}}s_W^2\overline{u}_i^c\gamma ^\mu u_i^c+{\displaystyle \frac{1}{3}}s_W^2\overline{d}_i^c\gamma ^\mu d_i^c`$ $`+(1s_W^2)(V_{i3}^uV_{3j}^u\overline{u}_i\gamma ^\mu u_j+V_{i3}^dV_{3j}^d\overline{d}_i\gamma ^\mu d_j)]`$ $`+{\displaystyle \frac{1}{2}}[\mathrm{\Delta }_{u_i}\overline{T}\gamma ^\mu u_i+\mathrm{\Delta }_{Dd_i}\overline{D}\gamma ^\mu d_i+\mathrm{\Delta }_{Sd_i}\overline{S}\gamma ^\mu d_i+\mathrm{h}.\mathrm{c}.]\},`$ where the leading-order coupling is given in terms of the standard fermion currents defined in Eq. (101). The $`Z`$ boson couples to pairs of heavy quarks at order one through the electromagnetic current $`J_Q`$. Note the flavor-changing couplings induced by $`ZZ^{}`$ mixing. The couplings of photons to fermions are given by the electromagnetic current as usual. #### B.2.5 Constraints from flavor physics: anomaly-free embedding The flavor-changing couplings of $`Z^{}`$ to quark pairs can feed into low-energy observables, leading to potentially large flavor-changing neutral currents. The contributions of the anomaly-free fermion embedding to mixing in the neutral $`K`$, $`D`$, $`B`$, and $`B_s`$ systems and the rare decays $`B_{d,s}\mu ^+\mu ^{}`$ and $`BK\mu ^+\mu ^{}`$ were summarized in Ref. in the context of 3-3-1 models without the little Higgs mechanism. If the quark mixing matrices take a Fritzsch-like structure , $`V_{ij}^{u,d}=\sqrt{m_j/m_i}`$ ($`ij`$), then the strongest bound on the $`Z^{}`$ mass comes from $`B`$$`\overline{B}`$ mixing and requires $`M_Z^{}>10.5`$ TeV . The next-most-stringent constraint comes from $`B_s`$$`\overline{B}_s`$ mixing and requires $`M_Z^{}>5.0`$ TeV. Clearly, the down quark mixing matrix must be more diagonal than the Fritzsch-like structure, in order to suppress flavor-changing effects in the down quark sector. In fact, one can choose $`V_{i3}^d=\delta _{i3}`$, so that the $`d`$ couplings are flavor-diagonal; this eliminates flavor-changing effects in the down quark sector. The flavor-changing effects are then pushed into the up sector. The $`u`$ and $`d`$ couplings to $`Z^{}`$ can never both be flavor-diagonal because they are related by the CKM matrix \[Eq. (112)\]. #### B.2.6 Quark masses and mixing: universal embedding In the universal embedding, the quark Yukawa Lagrangian is given for all three generations by $$=\lambda _1^{un}iu_1^{nc}\mathrm{\Phi }_1^{}Q_n+\lambda _2^{un}iu_2^{nc}\mathrm{\Phi }_2^{}Q_n+\frac{\lambda _d^{mn}}{\mathrm{\Lambda }}id_m^cϵ_{ijk}\mathrm{\Phi }_1^i\mathrm{\Phi }_2^jQ_n^k+\mathrm{h}.\mathrm{c}.,$$ (136) where $`m,n=1,2,3`$ are generation indexes; $`i,j,k=1,2,3`$ are SU(3) indexes; $`d_m^c`$ runs over all the down-type conjugate quarks ($`d^c,s^c,b^c`$); and $`u_{1,2}^{nc}`$ are linear combinations of the up-type conjugate quarks as given in Eqs. (139) and (141) below. The physics of the down quark sector in the universal embedding is exactly analogous to that of the charged leptons. The down quark Higgs couplings are given by $$=\frac{m_{d_i}}{v}y_d(Hd_i^cd_i)+\mathrm{h}.\mathrm{c}.,y_d=1\frac{v^2}{6f^2}(3+\frac{c_\beta ^4}{s_\beta ^2}+\frac{s_\beta ^4}{c_\beta ^2}),$$ (137) and their couplings to $`\eta `$ are given by $$=\frac{m_{d_i}}{v}\frac{v}{4\sqrt{2}f}\left(\frac{c_\beta ^2s_\beta ^2}{s_\beta c_\beta }\right)(i\eta d_i^cd_i)+\mathrm{h}.\mathrm{c}.$$ (138) In the up sector, the $`f`$ vevs generate mass terms for the three heavy quarks with charge $`+2/3`$. The three states $$U_n^c=\frac{\lambda _1^{un}c_\beta u_1^{nc}+\lambda _2^{un}s_\beta u_2^{nc}}{\sqrt{(\lambda _1^{un})^2c_\beta ^2+(\lambda _2^{un})^2s_\beta ^2}},$$ (139) marry the three $`U_n`$ states, giving them masses of $$M_{U_n}=f\sqrt{(\lambda _1^{un})^2c_\beta ^2+(\lambda _2^{un})^2s_\beta ^2}$$ (140) and leaving the orthogonal combinations of $`u_1^{nc}`$ and $`u_2^{nc}`$ massless: $$u_n^c=\frac{\lambda _2^{un}s_\beta u_1^{nc}+\lambda _1^{un}c_\beta u_2^{nc}}{\sqrt{(\lambda _1^{un})^2c_\beta ^2+(\lambda _2^{un})^2s_\beta ^2}}.$$ (141) Note that the Yukawa Lagrangian in Eq. (136) does not generate a misalignment between the SM up quark mass eigenstates and the heavy quarks. Such a misalignment could be generated by adding an additional dimension-5 operator, $$\frac{\lambda _u^{mn}}{\mathrm{\Lambda }}iu_m^cϵ_{ijk}\mathrm{\Phi }_1^i\mathrm{\Phi }_2^jQ_n^k+\mathrm{h}.\mathrm{c}.,$$ (142) to generate off-diagonal entries in the up quark mass matrix. We ignore this possibility here. The usual CKM matrix is generated by the off-diagonal entries in the down quark mass matrix, controlled by $`\lambda _d^{mn}`$. After EWSB, the up quark mass terms are $`_{\mathrm{up}\mathrm{mass}}`$ $`=`$ $`M_{U_n}U_n^cU_n+{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{s_\beta c_\beta [(\lambda _1^{un})^2(\lambda _2^{un})^2]}{\sqrt{(\lambda _1^{un})^2c_\beta ^2+(\lambda _2^{un})^2s_\beta ^2}}}U_n^cu_n`$ (143) $`{\displaystyle \frac{v}{\sqrt{2}}}{\displaystyle \frac{\lambda _1^{un}\lambda _2^{un}}{\sqrt{(\lambda _1^{un})^2c_\beta ^2+(\lambda _2^{un})^2s_\beta ^2}}}u_n^cu_n+\mathrm{h}.\mathrm{c}.`$ These terms lead to mixing between the heavy quarks and their corresponding SM quark partners. As usual, we use the subscript $`0`$ to denote SU(3) states; fields with no subscript denote the mass eigenstates after the mixing induced by EWSB. We can rewrite the SU(3) state $`U_{m0}`$ in terms of the mass eigenstate $`U_m`$ and the SM fermion $`u_m`$ as $$U_{m0}=U_m+\delta _{u_m}u_m,u_{m0}=\left(1\frac{1}{2}\delta _{u_m}^2\right)u_m\delta _{u_m}U_m,$$ (144) where $$\delta _{u_m}=\frac{v}{\sqrt{2}f}\frac{s_\beta c_\beta [(\lambda _1^{um})^2(\lambda _2^{um})^2]}{[(\lambda _1^{um})^2c_\beta ^2+(\lambda _2^{um})^2s_\beta ^2]}.$$ (145) The masses of the SM up-type quarks are given to leading order by $$\frac{\lambda _1^{um}\lambda _2^{um}f}{\sqrt{2}M_{U_m}}=\frac{m_{u_m}}{v}.$$ (146) The small mass of the $`u`$ ($`c`$) quark requires one of the couplings $`\lambda _{1,2}^{u1}`$ ($`\lambda _{1,2}^{u2}`$) to be very small. We choose the small coupling to be $`\lambda _1^{u1}`$ ($`\lambda _1^{u2}`$) so that the mixing effects in the up-quark sector are suppressed in the same $`t_\beta >1`$ limit as the mixing effects in the neutrino sector. We then have, $$M_U=f\lambda _Us_\beta ,M_C=f\lambda _Cs_\beta ,M_T=f\sqrt{\lambda _1^2c_\beta ^2+\lambda _2^2s_\beta ^2},$$ (147) where we define $`\lambda _U=\lambda _2^{u1}`$, $`\lambda _C=\lambda _2^{u2}`$, $`\lambda _1=\lambda _1^{u3}`$, and $`\lambda _2=\lambda _2^{u3}`$. For the mixing angles we also have $$\delta _u=\delta _c=\frac{v}{\sqrt{2}t_\beta f}=\delta _\nu ,\delta _t=\frac{vf}{\sqrt{2}M_T^2}s_\beta c_\beta (\lambda _1^2\lambda _2^2).$$ (148) We now write the up quark couplings to scalars. The couplings of heavy quark-partner pairs are given by $``$ $`=`$ $`(i\eta U^cU){\displaystyle \frac{c_\beta }{\sqrt{2}}}\lambda _U(i\eta C^cC){\displaystyle \frac{c_\beta }{\sqrt{2}}}\lambda _C+(i\eta T^cT)s_\beta c_\beta (\lambda _1^2\lambda _2^2){\displaystyle \frac{f}{\sqrt{2}M_T}}`$ (149) $`+(HT^cT){\displaystyle \frac{v}{f}}\left[(\lambda _1^2s_\beta ^2+\lambda _2^2c_\beta ^2){\displaystyle \frac{f}{2M_T}}s_\beta ^2c_\beta ^2(\lambda _1^2\lambda _2^2)^2{\displaystyle \frac{f^3}{2M_T^3}}\right]+\mathrm{h}.\mathrm{c}.`$ One would naively expect an $`HU_m^cU_m`$ coupling for the first two generations at order $`v/f`$ coming from replacing one Higgs field by its vev in the $`HHU_m^cU_m`$ term that is generated by the expansion of the nonlinear sigma model; however, this term is exactly canceled by the contribution from $`HU_m^cu_m`$ after $`uU`$ mixing in the first two generations if the up and charm quark masses are neglected. The leading-order couplings of the scalars to one heavy quark partner and one SM up-type quark are $$=(HU^cu)\frac{c_\beta \lambda _U}{\sqrt{2}}(HC^cc)\frac{c_\beta \lambda _C}{\sqrt{2}}+(HT^ct)(\lambda _1^2\lambda _2^2)\frac{s_\beta c_\beta f}{\sqrt{2}M_T}(i\eta t^cT)\frac{m_t}{v}+\mathrm{h}.\mathrm{c}.,$$ (150) where in the $`\eta `$ couplings we neglect $`m_u`$ and $`m_c`$ in the couplings of the first two generations and neglect the $`v/f`$ suppressed coupling of the third generation. These couplings will lead to the decays $`U_mu_mH`$ and $`Tt\eta `$. The couplings of scalars to a pair of SM up-type quarks are $``$ $`=`$ $`(Hu_i^cu_i)\left\{{\displaystyle \frac{m_{u_i}}{v}}\left[1{\displaystyle \frac{v^2}{6f^2}}\left({\displaystyle \frac{s_\beta ^4}{c_\beta ^2}}+{\displaystyle \frac{c_\beta ^4}{s_\beta ^2}}\right)\delta _{u_i}{\displaystyle \frac{v}{2\sqrt{2}f}}\left({\displaystyle \frac{c_\beta ^2s_\beta ^2}{s_\beta c_\beta }}\right)\right]+{\displaystyle \frac{M_{U_i}}{v}}\delta _{u_i}\delta _{u_i^c}\right\}`$ (151) $`+(i\eta u_i^cu_i){\displaystyle \frac{m_{u_i}}{v}}\left[{\displaystyle \frac{v}{\sqrt{2}f}}\left({\displaystyle \frac{s_\beta ^2c_\beta ^2}{s_\beta c_\beta }}\right)+\delta _{u_i}\right]+\mathrm{h}.\mathrm{c}.,`$ where the mixing between $`u_i^c`$ and $`U_i^c`$ at order $`v^2/f^2`$ is given by $`U_i^c=U_{i0}^c\delta _{u_i^c}u_{i0}^c`$, with $$\delta _{u_i^c}=\frac{m_{u_i}}{M_{U_i}}\left[\delta _{u_i}+\frac{v}{2\sqrt{2}f}\left(\frac{c_\beta ^2s_\beta ^2}{s_\beta c_\beta }\right)\right].$$ (152) #### B.2.7 Quark couplings to gauge bosons: universal embedding The couplings of the $`Z^{}`$ boson to quarks in the universal embedding were given in Table 4. These couplings are purely flavor-diagonal in the universal fermion embedding. The couplings of the heavy off-diagonal gauge bosons $`X^{}`$ and $`Y^0`$ to quarks in the universal embedding were also given in Table 4, neglecting CKM mixing. Keeping the full CKM dependence, we have $$_{X,Y}=\frac{g}{\sqrt{2}}[iX_\mu ^{}\overline{d}_i\gamma ^\mu (V_{ji}^{\mathrm{CKM}}U_j+\delta _{u_j}V_{ji}^{\mathrm{CKM}}u_j)+iY_\mu ^0\overline{u}_i\gamma ^\mu (U_i+\delta _{u_i}u_i)+\mathrm{h}.\mathrm{c}.].$$ (153) The couplings of $`W^\pm `$ to quark pairs, keeping terms of order $`v^2/f^2`$ in interactions involving only SM particles and terms of order $`v/f`$ in interactions involving one or more heavy particles, are $$_W=\frac{gW_\mu ^+}{\sqrt{2}}[(1\frac{1}{2}\delta _{u_i}^2)V_{ij}^{\mathrm{CKM}}\overline{u}_i\gamma ^\mu d_j\delta _{u_i}V_{ij}^{\mathrm{CKM}}\overline{U}_i\gamma ^\mu d_j+\mathrm{h}.\mathrm{c}.].$$ (154) The couplings of the $`Z`$ boson to quarks, including the corrections from mixing between $`Z`$ and $`Z^{}`$ and mixing between the TeV-scale quarks and their SM partners, are $`_Z`$ $`=`$ $`{\displaystyle \frac{gZ_\mu }{c_W}}\{(J_3^\mu s_W^2J_Q^\mu ){\displaystyle \frac{1}{2}}\delta _{u_i}^2\overline{u}_i\gamma ^\mu u_i{\displaystyle \frac{1}{2}}[\delta _{u_i}\overline{U}_i\gamma ^\mu u_i+\mathrm{h}.\mathrm{c}.]`$ $`+`$ $`{\displaystyle \frac{\delta _Z}{\sqrt{34s_W^2}}}[({\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{3}}s_W^2)(\overline{u}_i\gamma ^\mu u_i+\overline{d}_i\gamma ^\mu d_i){\displaystyle \frac{2}{3}}s_W^2\overline{u}_i^c\gamma ^\mu u_i^c+{\displaystyle \frac{1}{3}}s_W^2\overline{d}_i^c\gamma ^\mu d_i^c]\},`$ where the leading-order coupling is given in terms of the usual fermion currents $`J_3`$ and $`J_Q`$ defined in Eq. (101). The $`Z`$ boson couples to pairs of heavy quarks $`U_i`$ at order one through the electromagnetic current $`J_Q`$. The couplings of photons to fermions are given by the electromagnetic current as usual. ### B.3 Higgs potential In this section we describe the generation of the Higgs potential.<sup>9</sup><sup>9</sup>9We thank Martin Schmaltz for very helpful discussions. Additional details can be found in Refs. . We start with the Coleman-Weinberg potential that is generated by loops of gauge bosons and fermions in the running down from the cutoff scale $`\mathrm{\Lambda }`$. Above the global symmetry breaking scale $`f`$, only operators that are symmetric under the global \[SU(3)$`\times `$U(1)\]<sup>2</sup> symmetry are generated by the running. The three allowed operators up to dimension four are $$\mathrm{\Phi }_1^{}\mathrm{\Phi }_1,\mathrm{\Phi }_2^{}\mathrm{\Phi }_2,\|\mathrm{\Phi }_1^{}\mathrm{\Phi }_2\|^2.$$ (156) The first two of these operators are just constants and do not involve the Goldstone bosons. We therefore focus on the third operator. Expanding it in terms of the Goldstone bosons to fourth order gives $$\|\mathrm{\Phi }_1^{}\mathrm{\Phi }_2\|^2=f^4s_\beta ^2c_\beta ^2f^2h^{}h+\frac{1}{3s_\beta ^2c_\beta ^2}(h^{}h)^2+\frac{3}{32s_\beta ^2c_\beta ^2}h^{}h\eta ^2+𝒪(\varphi ^6).$$ (157) Running below the global symmetry breaking scale $`f`$ can give contributions to the Coleman-Weinberg potential that are not proportional to $`\|\mathrm{\Phi }_1^{}\mathrm{\Phi }_2\|^2`$. These contributions will contain logs of the ratios of masses-squared of $`f`$-scale particles and the corresponding SM particles. They will therefore be calculable, i.e., independent of cutoff-scale physics. The Coleman-Weinberg potential from the $`X^{}`$, $`Y^0`$ and $`W^+`$ gauge bosons is, $`V_2`$ $`=`$ $`{\displaystyle \frac{3}{64\pi ^2}}g^4\mathrm{log}(\mathrm{\Lambda }^2/M_X^2)f^2(h^{}h)`$ $`V_4`$ $`=`$ $`{\displaystyle \frac{3}{64\pi ^2}}g^4\mathrm{log}(\mathrm{\Lambda }^2/M_X^2)\left[{\displaystyle \frac{1}{3s_\beta ^2c_\beta ^2}}(h^{}h)^2{\displaystyle \frac{3}{32s_\beta ^2c_\beta ^2}}(h^{}h)\eta ^2\right]`$ (158) $`{\displaystyle \frac{3}{128\pi ^2}}g^4\mathrm{log}(M_X^2/M_W^2)(h^{}h)^2.`$ Here $`V_2`$ and the first line of $`V_4`$ come from running between $`\mathrm{\Lambda }`$ and $`M_X`$ and are proportional to $`\|\mathrm{\Phi }_1^{}\mathrm{\Phi }_2\|^2`$, while the second line of $`V_4`$ comes from running between $`M_X`$ and $`M_W`$. The running below $`M_X`$ contributes only a term involving $`(h^{}h)^2`$. It does not contribute any terms involving $`\eta `$ since there is no coupling of $`W`$ boson pairs to $`h\eta `$. The Coleman-Weinberg potential from the $`Z^{}`$ and $`Z`$ gauge bosons is, $`V_2`$ $`=`$ $`{\displaystyle \frac{3}{32\pi ^2}}g^4{\displaystyle \frac{1+t_W^2}{3t_W^2}}\mathrm{log}(\mathrm{\Lambda }^2/M_Z^{}^2)f^2(h^{}h)`$ $`V_4`$ $`=`$ $`{\displaystyle \frac{3}{32\pi ^2}}g^4{\displaystyle \frac{1+t_W^2}{3t_W^2}}\mathrm{log}(\mathrm{\Lambda }^2/M_Z^{}^2)\left[{\displaystyle \frac{1}{3s_\beta ^2c_\beta ^2}}(h^{}h)^2{\displaystyle \frac{3}{32s_\beta ^2c_\beta ^2}}(h^{}h)\eta ^2\right]`$ (159) $`{\displaystyle \frac{3}{256\pi ^2}}g^4(1+t_W^2)^2\mathrm{log}(M_Z^{}^2/M_Z^2)(h^{}h)^2.`$ Again, $`V_2`$ and the first line of $`V_4`$ come from running between $`\mathrm{\Lambda }`$ and $`M_Z^{}`$ and are proportional to $`\|\mathrm{\Phi }_1^{}\mathrm{\Phi }_2\|^2`$, while the second line of $`V_4`$ comes from running between $`M_Z^{}`$ and $`M_Z`$. The running below $`M_Z^{}`$ contributes only a term involving $`(h^{}h)^2`$. It does not contribute any terms involving $`\eta `$ since there is no coupling of $`Z`$ boson pairs to $`h\eta `$. The Coleman-Weinberg potential from the fermions can in principle come from loops of any fermion with an order-one Yukawa coupling. However, due to the feature of collective breaking in the model, the order-one Yukawa couplings that give mass to the neutrino partners and the quark partners of the first two generations do not contribute to the terms of the Coleman-Weinberg potential involving the Goldstone bosons (neglecting the tiny Yukawa couplings of the quarks of the first two generations). The only significant contribution is then due to the top quark and its partner $`T`$. In what follows we neglect the mixing between quark generations. The Coleman-Weinberg potential from the top quark and its partner $`T`$ is, $`V_2`$ $`=`$ $`{\displaystyle \frac{3}{8\pi ^2}}\lambda _t^2M_T^2\mathrm{log}(\mathrm{\Lambda }^2/M_T^2)(h^{}h)`$ $`V_4`$ $`=`$ $`{\displaystyle \frac{3}{8\pi ^2}}\lambda _t^2{\displaystyle \frac{M_T^2}{f^2}}\mathrm{log}(\mathrm{\Lambda }^2/M_T^2)\left[{\displaystyle \frac{1}{3s_\beta ^2c_\beta ^2}}(h^{}h)^2{\displaystyle \frac{3}{32s_\beta ^2c_\beta ^2}}(h^{}h)\eta ^2\right]`$ (160) $`+{\displaystyle \frac{3}{16\pi ^2}}\lambda _t^4\mathrm{log}(M_T^2/m_t^2)(h^{}h)^2,`$ where $`\lambda _t\lambda _1^t\lambda _2^tf/M_T\sqrt{2}m_t/v`$. Again, $`V_2`$ and the first line of $`V_4`$ come from running between $`\mathrm{\Lambda }`$ and $`M_T`$ and are proportional to $`\|\mathrm{\Phi }_1^{}\mathrm{\Phi }_2\|^2`$, while the second line of $`V_4`$ comes from running between $`M_T`$ and $`m_t`$. The running below $`M_T`$ contributes only a term involving $`(h^{}h)^2`$. It does not contribute any terms involving $`\eta `$ since there is no coupling of top quark pairs to $`h\eta `$ or $`\eta ^2`$. Collecting terms, we can write the Coleman-Weinberg potential as follows: $$V=m^2h^{}h+\lambda (h^{}h)^2+\lambda ^{}h^{}h\eta ^2,$$ (161) where $`m^2`$ $`=`$ $`{\displaystyle \frac{3}{8\pi ^2}}\left[\lambda _t^2M_T^2\mathrm{log}(\mathrm{\Lambda }^2/M_T^2){\displaystyle \frac{g^2}{4}}M_X^2\mathrm{log}(\mathrm{\Lambda }^2/M_X^2){\displaystyle \frac{g^2}{8}}(1+t_W^2)M_Z^{}^2\mathrm{log}(\mathrm{\Lambda }^2/M_Z^{}^2)\right]`$ $`\lambda `$ $`=`$ $`{\displaystyle \frac{1}{3s_\beta ^2c_\beta ^2}}{\displaystyle \frac{m^2}{f^2}}+{\displaystyle \frac{3}{16\pi ^2}}\left[\lambda _t^4\mathrm{log}(M_T^2/m_t^2){\displaystyle \frac{g^4}{8}}\mathrm{log}(M_X^2/M_W^2){\displaystyle \frac{g^4}{16}}(1+t_W^2)^2\mathrm{log}(M_Z^{}^2/M_Z^2)\right]`$ $`\lambda ^{}`$ $`=`$ $`{\displaystyle \frac{3}{32s_\beta ^2c_\beta ^2}}{\displaystyle \frac{m^2}{f^2}}.`$ (162) In the expression for $`m^2`$, in principle the cutoff $`\mathrm{\Lambda }`$ in the term generated by quark loops can be different from the cutoff $`\mathrm{\Lambda }`$ in the two terms generated by gauge boson loops, because the physics that cuts off the quark loops can be different from the physics that cuts off the gauge boson loops. After EWSB, $`\eta `$ gets a small positive mass-squared of order $`m_H^2v^2/f^2`$ from the $`\lambda ^{}`$ term. The Higgs vev and mass are given by $$v^2=m^2/\lambda =(246\mathrm{GeV})^2,m_H^2=2m^2=2\lambda v^2.$$ (163) It turns out that this $`m_H`$ is too small, because the quartic coupling $`\lambda `$ is not big enough compared to $`m^2`$. Following Ref. , this problem can be fixed by adding a new operator, $`\mathrm{\Phi }_1^{}\mathrm{\Phi }_2+\mathrm{h}.\mathrm{c}.`$, to the scalar potential with a coefficient $`\mu ^2`$ set by hand. This operator breaks the global SU(3)<sup>2</sup> down to the diagonal SU(3) while preserving the gauged SU(3). Expanding this operator to fourth order in the Goldstone bosons gives $$\mathrm{\Phi }_1^{}\mathrm{\Phi }_2+\mathrm{h}.\mathrm{c}.=2f^2s_\beta c_\beta +\frac{1}{f^2s_\beta c_\beta }[f^2(h^{}h)\frac{f^2\eta ^2}{2}+\frac{(h^{}h)^2}{12s_\beta ^2c_\beta ^2}+\frac{3(h^{}h)\eta ^2}{32s_\beta ^2c_\beta ^2}+\frac{\eta ^4}{48s_\beta ^2c_\beta ^2}].$$ (164) Because the $`(h^{}h)`$ and $`(h^{}h)^2`$ terms in this operator have different relative coefficients than in the original operator $`\|\mathrm{\Phi }_1^{}\mathrm{\Phi }_2\|^2`$, it can be used to cancel off part of the $`m^2h^{}h`$ term without canceling too much of the $`\lambda (h^{}h)^2`$ term. Adding the term $`\mu ^2(\mathrm{\Phi }_1^{}\mathrm{\Phi }_2+\mathrm{h}.\mathrm{c}.)`$ to the potential gives $$V=m_{\mathrm{new}}^2h^{}h+\frac{1}{2}m_\eta ^2\eta ^2+\lambda _{\mathrm{new}}(h^{}h)^2+\lambda _{\mathrm{new}}^{}h^{}h\eta ^2+\lambda _{\mathrm{new}}^{\prime \prime }\eta ^4,$$ (165) where $`m_{\mathrm{new}}^2`$ $`=`$ $`m^2{\displaystyle \frac{\mu ^2}{s_\beta c_\beta }},m_\eta ^2={\displaystyle \frac{\mu ^2}{s_\beta c_\beta }},`$ $`\lambda _{\mathrm{new}}`$ $`=`$ $`{\displaystyle \frac{1}{3s_\beta ^2c_\beta ^2}}{\displaystyle \frac{m_{\mathrm{new}}^2}{f^2}}+{\displaystyle \frac{1}{4s_\beta ^3c_\beta ^3}}{\displaystyle \frac{\mu ^2}{f^2}}`$ $`+{\displaystyle \frac{3}{16\pi ^2}}\left[\lambda _t^4\mathrm{log}(M_T^2/m_t^2){\displaystyle \frac{g^4}{8}}\mathrm{log}(m_X^2/m_W^2){\displaystyle \frac{g^4}{16}}(1+t_W^2)^2\mathrm{log}(m_Z^{}^2/m_Z^2)\right],`$ $`\lambda _{\mathrm{new}}^{}`$ $`=`$ $`{\displaystyle \frac{3}{32s_\beta ^2c_\beta ^2}}{\displaystyle \frac{m_{\mathrm{new}}^2}{f^2}},\lambda _{\mathrm{new}}^{\prime \prime }={\displaystyle \frac{1}{48s_\beta ^3c_\beta ^3}}{\displaystyle \frac{\mu ^2}{f^2}}.`$ (166) Note that this term has also given rise to a mass-squared term for $`\eta `$ and an $`\eta ^4`$ coupling. The $`\eta `$ mass $`m_\eta `$ is now of order $`\mu `$, parametrically larger than the $`\eta `$ mass term generated by EWSB.
warning/0506/hep-th0506022.html	ar5iv	text	# Contents ## 1 Introduction The recent revival of interest in cosmic strings is due to developments on both observational and theoretical fronts. On the observational side there is optimism that the next generation of gravitational wave detectors, Advanced LIGO and LISA, will be able to detect the characteristic signature of cosmic string loops as they twist and whip into cusps . Moreover, there exist tantalizing gravitational lensing events, most notably CSL1 and its companion images , which point at the existence of a cosmic string with tension $`G\mu 4\times 10^7`$. Further spectroscopic analysis should reveal the nature of this system in the near future. Reviews of these developments can be found in . On the theoretical side, the advent of the “warped throat” in realistic type II string compactifications has permitted a resurrection of the (b)old idea that cosmic strings may be superstrings stretched across the sky . In this modern guise, the cosmic string network may consist of both fundamental strings, D-strings and wrapped D-branes. Reviews of these recent stringy developments can be found in and , while earlier work on cosmic strings from the 1980’s and early 90’s is summarized in . Having admitted the theoretical possibility that cosmic strings may be fundamental strings, the important question becomes: how can we tell? As described in , there are two major distinctions between fundamental strings and gauge theoretic solitons in perturbative field theories. The first is the existence of multi-tension string networks consisting of both fundamental strings, and D-strings, and their various bound states. The study of the dynamical scaling properties of such networks is underway . The second distinguishing feature of fundamental strings, and the one we focus on in this paper, is their interaction cross-section. In the abelian-Higgs model it is known that cosmic strings reconnect with unit probability $`P=1`$ over a wide range of impact parameters. In contrast, fundamental strings interact with probability $`Pg_s^2`$, where the functional dependence on the angle of incidence and relative velocity of the strings was determined in . Similarly, D-strings may also pass through each other, reconnecting with probability $`P<1`$ . A reduced probability for reconnection affects the scaling solution for the string network, resulting in a larger concentration of strings in the sky . Given enough data, it is not implausible that one could extract the probability $`P`$ from observation. For further work on cosmic superstrings, see . Of course, it may be possible to engineer a gauge theory whose solitonic strings mimic the behavior of fundamental strings. Indeed, since strongly coupled gauge theories are often dual to string theories in warped throats, this must be true on some level. However, restricting attention to weakly coupled gauge theories, we could ask if there exist semi-classical cosmic strings which, like their fundamental cousins, reconnect with probability $`P<1`$. A mechanism for achieving this was mentioned by Polchinski in : construct a vortex with extra internal bosonic zero modes. Two vortices could then miss each other in this internal space. It was further noted that symmetry breaking effects would generically give mass to these internal modes, ruining the mechanism except in rather contrived models. In this paper we present a model which realizes Polchinski’s field theoretic counterexample although, ultimately, in a rather different manner than anticipated. Our model embeds the cosmic string in a non-abelian $`U(N)`$ gauge theory, so that the string may move in the internal color and flavor space. The formal properties of vortices of this type have been studied extensively of late (see ) although, until now, not in the context of cosmic strings. Despite the presence of this internal space we find, rather surprisingly, that cosmic strings continue to reconnect with essentially unit probability, passing through each other only for finely tuned initial conditions. This result occurs due to the non-trivial topology in the interior region of the two-vortex moduli space. However, we show that this conclusion is changed when the internal modes gain a mass. Lifting the vortex zero modes naturally leaves behind $`N`$ different cosmic strings, each of which reconnects only with strings of the same type while passing through other types of strings. The effect of lifting the internal moduli space is therefore to reduce the probability of reconnection at low-energies to $`1/N`$. At high energies these effects wash-out, the strings may evolve into each other, and reconnection again occurs with probability one. In the next section we review the cosmic strings of interest and explain how they gain an internal space of massless modes. In section 3 we study the moduli space of two vortex strings and argue that vortex strings only fail to reconnect for a set of initial conditions of measure zero. In section 4 we examine various further effects in the model, including quantum dynamics on the vortex worldvolume, symmetry breaking masses and fermionic zero modes, and show how these effects reduce the probability of reconnection at low energies. We conclude with the traditional conclusions. ## 2 Non-Abelian Cosmic Strings In this paper we study the dynamics of cosmic strings living in a non-abelian $`U(N_c)`$ gauge theory, coupled to $`N_f`$ scalars $`q_i`$, each transforming in the fundamental representation, $`L={\displaystyle \frac{1}{4e^2}}\mathrm{Tr}F_{\mu \nu }F^{\mu \nu }+{\displaystyle \underset{i=1}{\overset{N_f}{}}}𝒟_\mu q_i^{}𝒟_\mu q_i{\displaystyle \frac{\lambda e^2}{2}}\mathrm{Tr}\left({\displaystyle \underset{i=1}{\overset{N_f}{}}}q_iq_i^{}v^2\right)^2`$ (2.1) When $`N_c=N_f=1`$, this is simply the abelian-Higgs model while, for $`N_c>1`$, it is the simplest non-abelian generalization. The Lagrangian (2.1) enjoys a $`SU(N_f)`$ flavor symmetry, rotating the scalars. In Section 4 we shall discuss the more realistic situation in which this symmetry is softly broken, but for now let us assume it remains intact. Moreover, we shall restrict attention to the case $`N_f=N_cN`$<sup>1</sup><sup>1</sup>1When $`N_f<N_c`$, the central $`U(1)`$ remains unbroken and the theory does not admit vortices. For $`N_f>N_c`$, the resulting cosmic strings are non-abelian generalizations of semi-local vortices .. The theory (2.1) has a unique vacuum in which the scalars condense in the pattern $`q_i^a=v\delta _i^a`$ (2.2) Here $`i=1,\mathrm{},N`$ is the flavor index, while $`a=1,\mathrm{}N`$ is the color index. In what follows we will take our theory to be weakly coupled by requiring the symmetry breaking scale $`ev\mathrm{\Lambda }`$, the scale at which the non-abelian sector confines. The vacuum (2.2) has a mass gap in which the gauge field has mass $`m_\gamma ^2e^2v^2`$, while the the scalars have mass $`m_q^2\lambda e^2v^2`$. The symmetry breaking pattern resulting from this condensate puts the theory into what is referred to as the “color-flavor-locked” phase, with $`U(N_c)\times SU(N_f)SU(N)_{\mathrm{diag}}`$ (2.3) The fact that the overall $`U(1)U(N_c)`$ is broken in the vacuum guarantees the existence of vortex strings characterized by the winding $`\mathrm{Tr}B=2\pi k`$ for some $`k𝐙`$, where the integral is over the plane transverse to the vortex. The tension of this cosmic string is given by $`T=2\pi v^2f(\lambda )`$ (2.4) where $`f(\lambda )`$ is a slowly varying function with $`f(1)=1`$. The width of the vortex core is given by $`L\mathrm{max}(1/m_\gamma ,1/m_q)`$. The parameter $`\lambda `$ dictates the behavior of multiple, parallel vortex strings: when $`\lambda >1`$, parallel vortex strings repel (as in a type II superconductor) while, for $`\lambda <1`$, parallel strings attract (type I superconductor). In both cases, the force is short ranged, dying off exponentially away from the vortex core. For the critical coupling $`\lambda =1`$, vortex strings feel no force and multi-soliton solutions exist with parallel vortex strings sitting at arbitrary positions. Vortex solutions to the theory (2.1) can easily be constructed from solutions to the corresponding abelian theory. Let $`A_\mu ^{}`$ and $`q^{}`$ denote gauge and Higgs profiles of the abelian vortex solution. Then we can construct a non-abelian vortex by simply embedding in the upper-left-hand corner thus: $`q_i^a=\left(\begin{array}{cccc}q^{}& & & \\ & v& & \\ & & \mathrm{}& \\ & & & v\end{array}\right)(A_\mu )_b^a=\left(\begin{array}{cccc}A_\mu ^{}& & & \\ & 0& & \\ & & \mathrm{}& \\ & & & 0\end{array}\right)`$ (2.13) This is not the most general embedding. We can act on this configuration with the $`SU(N)_{\mathrm{diag}}`$ symmetry preserved by the vacuum to generate new solutions. Dividing out by the stabilizing group, the space of vortex solutions is $`{\displaystyle \frac{SU(N)_{\mathrm{diag}}}{SU(N1)\times U(1)}}^{N1}`$ (2.14) The existence of these internal, Goldstone modes, on the string worldsheet means that, at low-energies, the string feels as if it is propagating in a higher dimensional space $`𝐑^{3,1}\times ^{N1}`$. The size (Kähler class) of the internal $`^{N1}`$ space is given by $`r={\displaystyle \frac{\stackrel{~}{f}(\lambda )}{e^2}}`$ (2.15) where $`\stackrel{~}{f}(\lambda )`$ is, once again, a slowly varying function of $`\lambda `$ and it can be shown that $`\stackrel{~}{f}(1)=2\pi `$. A few comments on the literature: vortex zero modes arising through a mechanism of this type were previously studied in although these authors considered unbroken gauge symmetries, a situation which leads to further subtleties. The term “non-abelian strings” is also used to refer to simply-connected gauge groups broken to a discrete subgroup, often giving rise to several types of cosmic string; see for example . Such strings have rather different properties from those considered here, such as the existence of string junctions, and their dynamics shares features with $`(p,q)`$ string networks . Finally, we make no attempt to embed our model in a viable GUT, preferring to concentrate instead on the robust features of our vortex strings. A detailed description of $`SO(10)`$ GUT strings can be found, for example, in . The cosmological consequences of these internal modes mimic the behavior of string moving in higher dimensions. The internal currents carried by the string are akin to motion in the higher dimensions and, through equipartition of energy, have the effect of slowing down the motion of the strings in the three dimensions of real space . As we will explain in Section 4, in our case these internal modes actually gain a small mass from quantum effects and such currents cease to play a role over large times. More important for the present discussion is the fact that the internal space (2.14), like the higher dimensions of string theory, allows vortices to pass without interacting. To see this, consider two abelian vortices in orthogonal $`U(1)`$ subgroups, but at different points $`x_1`$ and $`x_2`$ in space, $`q_i^a=\left(\begin{array}{cccc}q^{}(x_1)& & & \\ & q^{}(x_2)& & \\ & & \mathrm{}& \\ & & & v\end{array}\right)(A_\mu )_b^a=\left(\begin{array}{cccc}A_\mu ^{}(x_1)& & & \\ & A_\mu ^{}(x_2)& & \\ & & \mathrm{}& \\ & & & 0\end{array}\right)`$ (2.24) In this case, the two strings evolve independently and simply pass through each other: no reconnection occurs. Of course, if the two vortices instead lie in the same $`U(1)`$ subgroup, as in (2.13), then vortices strongly interact and, as we review in the next section, reconnect. The question is: what happens in the intermediate situations when the two vortices lie in overlapping $`U(1)`$ subgroups? For fundamental strings moving in $`d`$ compactified dimensions of string theory, one expects the classical probability of reconnection to be suppressed by the geometric factor $`l_s^d/V`$, where $`V`$ is the volume of the extra dimensions and $`l_s=\sqrt{\alpha ^{}}`$ is the width of the string . Together with the inherent probability $`Pg_s^2`$ of fundamental string reconnection <sup>2</sup><sup>2</sup>2In realistic string compactifications, potentials on the internal space suppress the geometrical suppression while the $`g_s`$ suppression remains ., this leads to a reduced string cross-section, the net result of which is to increase the number of cosmic strings seen in the sky . Naively one may imagine that our field theoretic model exhibits similar behavior, with a critical separation in the internal space distinguishing reconnecting strings from those that pass through each other. This would then allow one to define a classical probability $`P`$ of reconnection in this system by coarse graining over the internal space. In the next section we turn to a detailed study of this issue. We shall find that the vortices always reconnect unless they lie in orthogonal subgroups. In some sense, the situation of orthogonal subgroups (2.24) is already the critical separation and the classical probability of reconnection is $`P=1`$. ## 3 The Classical Reconnection of Cosmic Strings In general the non-linear evolution of solitons is a difficult question that requires numerical investigation. However, for the low-energy scattering of cosmic strings we may reliably employ analytical methods in which we restrict attention to the light degrees of freedom describing the positions and internal orientations of the two strings. This method, known as the moduli space approximation , has been successfully applied to the abelian Higgs model where it was used to show that vortex strings indeed reconnect . Later numerical simulations revealed that this result is robust, holding for very high energy collisions . These results underpin the statement that gauge theoretic cosmic strings reconnect with probability one. Here we present the moduli space analysis for the non-abelian strings; it is to be hoped that a similar robustness holds for the present result. ### 3.1 Reconnection of $`U(1)`$ Strings Let us start by recalling how we see reconnection from the moduli space perspective in the case of the abelian Higgs model . One can reduce the dynamics of cosmic strings to that of particles by considering one of two spatial slices shown in Figure 1. The vertical slice cuts the strings to reveal a vortex-anti-vortex pair. After reconnection, this slice no longer intersects the strings, implying the annihilation of this pair. Alternatively, one can slice horizontally to reveal two vortices. Here the smoking gun for reconnection is the right-angle scattering of the vortices at (or near) the interaction point, as shown in Figure 1 (right). Such $`90^\mathrm{o}`$ degree scattering is a requirement since, as is clear from the figure, the two ends of each string are travelling in opposite directions after the collision. By varying the slicing along the string, one can reconstruct the entire dynamics of the two strings in this manner and show the inevitability of reconnection at low-energies. Hence, reconnection of cosmic strings requires both the annihilation of vortex-anti-vortex pairs and the right-angle scattering of two vortices. While the former is expected (at least for suitably slow collisions), to see the latter we must take a closer look at the dynamics of vortices. At critical coupling $`\lambda =1`$, the static forces between vortices cancel and we may rigorously define the moduli space of solutions to the vortex equations. The relative moduli space of two abelian vortices is simply $`𝐂/𝐙_2`$, where $`𝐂`$ is parameterized by $`z`$, the separation between vortices, and $`𝐙_2:zz`$ reflects the fact that the vortices are indistinguishable objects. This $`𝐙_2`$ action means that the single valued coordinate on the moduli space is $`z^2`$, rather than $`z`$, an important point in what follows. While the metric on this space is unknown<sup>1</sup><sup>1</sup>1Various properties of the metric on the vortex moduli space have been uncovered in ., it is known to be smooth , looking like the snub-nose cone shown in Figure 2. The motion of two particles at zero impact parameter goes up and over the cone, as shown in the figure, returning down the other side. This motion doesn’t correspond to scattering by $`180^\mathrm{o}`$ (this would be coming back down the same side), but to $`90^\mathrm{o}`$ scattering. This result does not depend on details of the metric on the vortex moduli space, but follows simply from the fact that, near the origin, the space is smooth and the single valued coordinate is $`z^2`$, rather than $`z`$. Before proceeding, we pass some well-known comments on the validity of the moduli space approximation. The energies involved in the collision should be small enough so as not to excite radiation. In the present context, this means that the relative velocity $`\dot{z}`$ of the vortices should satisfy $`TL\dot{z}^2m_\gamma ,m_q`$ where $`L\mathrm{max}(1/m_\gamma ,1/m_q)`$ is the width of the vortex string. For $`\lambda 1`$ this translates into the requirement that $`\dot{z}^2e^2`$. Similarly, the angle of incidence of the vortices, measured by $`z^{}`$, the spatial derivative of the separation along the string, should also satisfy $`z^2e^2`$. Finally, we should mention that the description in terms of particle motion on the moduli space is a little misleading, since a given slice of the string need not follow a geodesic on the moduli space. (For example, waves propagating along the string do not have this property). One should talk instead in terms of the dynamics of the real line $`𝐑`$, the spatial extent of the two strings, mapped into the moduli space. The inevitability of reconnection then follows from the single valued nature of $`z^2`$, together with the free motion of the strings far from the interaction point . ### 3.2 The Moduli Space of Non-Abelian Vortices We would now like to repeat this analysis for the the non-abelian vortices introduced in Section 2. For the vertical slice shown in Figure 1, the abelian argument carries over. Our vortices have only a single topological protector, $`k=\mathrm{Tr}B/2\pi `$, and a vortex-anti-vortex pair may annihilate regardless of their mutual orientation in the gauge group. One caveat is that a vortex and anti-vortex in orthogonal $`U(1)`$’s (as in (2.24)) remains a solution, albeit an unstable one. Therefore if we do not allow fluctuations away from this ansatz, the vortex-anti-vortex pair will pass right through each other without annihilating, in accord with the statements in the previous section. This is our first hint that reconnection will occur except for finely tuned initial conditions. To complete the argument of reconnection, we also need to study when right-angle scattering occurs. For this we require a description of the moduli space of multiple non-abelian vortices in the critically coupled ($`\lambda =1)`$ theory (2.1). A description of this space arising from modelling the system in terms of string theoretic D-branes was presented in . The moduli space is presented in terms of an algebraic quotient construction, related to the ADHM construction of the instanton moduli space. We now review this construction. The moduli space of $`k`$ vortices in $`U(N)`$ gauge theory is a Kähler manifold with real dimension $`2kN`$ which we denote as $`_{k,N}`$. The construction of presents this space as a $`U(k)`$ symplectic quotient construction. We start with a $`k\times k`$ complex matrix $`Z`$, and a $`k\times N`$ complex matrix $`\mathrm{\Psi }`$, subject to the constraint $`[Z,Z^{}]+\mathrm{\Psi }\mathrm{\Psi }^{}=r`$ (3.1) where the right-hand-side is proportional to the $`k\times k`$ identity matrix, and $`r=2\pi /e^2`$ as in (2.15). The moduli space $`_{k,N}`$ is defined as the quotient of the solutions to this constraint, where we divide by the $`U(k)`$ action $`ZUZU^{},\mathrm{\Psi }U\mathrm{\Psi }`$ (3.2) The $`U(k)`$ action has no fixed points ensuring that, as in the abelian case, the moduli space of vortices is smooth. Roughly speaking, the eigenvalues of $`Z`$ correspond to the positions<sup>1</sup><sup>1</sup>1The eigenvalues of $`Z`$ are dimensionless and correspond to the positions of the vortices multiplied by the mass scale $`v`$. of $`k`$ vortices, while the independent components of $`\mathrm{\Psi }`$ denote the orientations of these vortices in the internal space. For example, when $`k=1`$, the scalar $`Z`$ decouples and corresponds to the center of mass of the vortex, while $`\mathrm{\Psi }`$ satisfies $`\|\mathrm{\Psi }\|^2=r`$, modulo the $`U(1)`$ gauge action, which reproduces the moduli space $`^{N1}`$ for a single vortex. The manifold $`_{k,N}`$ has an $`SU(N)_{\mathrm{diag}}\times U(1)_R`$ action. The former results from the symmetry (2.3) acting on the vortex; the latter is the rotational symmetry of the plane. In the construction described above, the action is $`SU(N)_{\mathrm{diag}}:\mathrm{\Psi }\mathrm{\Psi }V,U(1)_R:Ze^{i\alpha }Z`$ (3.3) The algebraic quotient description of the vortex moduli space presented here arises from studying vortices in a Hanany-Witten set-up . We stress that, despite the D-brane origin of this construction, the resulting moduli space is that of field theoretic vortices; indeed, in , this framework was used to elucidate the differences between abelian ($`N=1`$) vortex strings and D-strings moving in vacua. Here we are interested in non-abelian ($`N2`$) strings. To our knowledge, there is no field theoretic derivation that the quotient space $`_{k,N}`$ coincides with the vortex moduli space and such a proof would be desirable. In the following we shall see that several key features of $`_{k,N}`$ correctly capture the behavior of vortex strings. Note however that the space $`_{k,N}`$ naturally inherits a metric from the above construction; this does not coincide with the metric on the moduli space of vortices (interpreted in terms of solitons, it describes co-dimension two objects with long-range polynomial tails). Thankfully, in what follows we will only require topological information about $`_{k,N}`$. ### 3.3 Reconnection of $`U(2)`$ Strings We first discuss the case of $`k=2`$ vortices in the $`N_c=N_f=2`$ gauge theory. Both $`Z`$ and $`\mathrm{\Psi }`$ are $`2\times 2`$ matrices (although for different reasons), and each suffers a $`U(2)`$ action (3.2). We project out the trivial center of mass motion of the system by requiring $`\mathrm{Tr}Z=0`$ and, following , use the $`U(2)`$ action to place $`Z`$ in upper-triangular form. We write $`Z=\left(\begin{array}{cc}z& \omega \\ 0& z\end{array}\right),\mathrm{\Psi }=\left(\begin{array}{cc}a_1& a_2\\ b_1& b_2\end{array}\right)`$ (3.8) This does not fix all gauge degrees of freedom, but leaves a surviving $`U(1)_1\times U(1)_2\times 𝐙_2U(2)`$ gauge symmetry acting as: $`U(1)_1:U=\left(\begin{array}{cc}e^{i\varphi }& 0\\ 0& 0\end{array}\right),U(1)_2:U=\left(\begin{array}{cc}0& 0\\ 0& e^{i\varphi }\end{array}\right)`$ (3.13) under which $`a_i`$ transforms with charge $`(1,0)`$, $`b_i`$ with charge $`(0,1)`$ and $`\omega `$ with charge $`(1,1)`$. The coordinate $`z`$ is neutral. Meanwhile, the $`𝐙_2`$ action is $`𝐙_2:U={\displaystyle \frac{1}{\sqrt{1+\|\zeta \|^2}}}\left(\begin{array}{cc}1& \overline{\zeta }\\ \zeta & 1\end{array}\right)`$ (3.16) with the parameter $`\zeta =2z/\omega `$. Finally, the constraints arising from (3.1) read $`{\displaystyle \underset{i=1}{\overset{2}{}}}\|a_i\|^2=r\|\omega \|^2,{\displaystyle \underset{i=1}{\overset{2}{}}}\|b_i\|^2=r+\|\omega \|^2,a_1\overline{b}_1+a_2\overline{b}_2=2\overline{z}\omega `$ (3.17) Counting degrees of freedom, we have 6 complex parameters in $`Z`$ and $`\mathrm{\Psi }`$ subject to two real constraints and one complex constraint (3.17), together with the two $`U(1)`$ actions (3.13). This leaves us with a smooth moduli space $`\stackrel{~}{}_{2,2}`$ of 3 complex dimensions. (Recall that we have factored out the center of mass degree of freedom so the full moduli space is $`_{2,2}𝐂\times \stackrel{~}{}_{2,2}`$). The rest of this section is devoted to studying this space. #### The Asymptotic Regime To get a feel for the physical interpretation of the various parameters, it is instructive to examine the regime of far separated vortices with $`zr`$. We have $`\|\omega \|1/\|z\|`$ and the constraints (3.17), combined with the $`U(1)^2`$ action (3.13), restrict $`a_i`$ and $`b_i`$ to lie in independent $`^1`$’s, up to $`1/\|z\|^2`$ corrections. In this limit the $`𝐙_2`$ action reads $`𝐙_2:\{\begin{array}{ccc}z& & z\\ a_i& & b_i\end{array}`$ (3.20) interchanging the position and orientation of the two vortices. Thus, asymptotically, the moduli space is simply $`\stackrel{~}{}_{2,2}{\displaystyle \frac{𝐂\times ^1\times ^1}{𝐙_2}}`$ (3.21) #### Two Submanifolds To continue our exploration of this space, it will prove useful to seek out a couple of special submanifolds. These correspond to the two extreme cases described in Section 2 in which we understand that reconnection does/does not occur. The first such submanifold corresponds to the situation (2.13) where the vortices lie in the same $`U(1)`$ subgroup. As we mentioned in Section 2, such vortices must always scatter at right-angles. We can impose this condition through the requirement that the two orientation vectors lie parallel: $`a_ib_i`$. We will refer to this submanifold as $`\|_{U(1)}\stackrel{~}{}_{2,2}`$. By an $`SU(2)_F`$ action, we can choose a representative point, say $`a_2=b_2=0`$. Then the constraints read $`\|a_1\|^2=r\|\omega \|^2,\|b_1\|^2=r+\|\omega \|^2,a_1\overline{b}_1=2\overline{z}\omega `$ (3.22) This system was previously studied in . On this submanifold, $`a_i`$ and $`b_i`$ are both even under the $`𝐙_2`$ action (3.16), while $`\omega `$ and $`z`$ are odd: $`(\omega ,z)(\omega ,z)`$. The calculation of shows that this manifold is asymptotic to $`𝐂/𝐙_2`$, with a smooth metric at the origin, as depicted in Figure 2. Acting now with the $`SU(2)_F`$ action sweeps out a $`^1`$ at each point, leaving us with a space which is topologically<sup>1</sup><sup>1</sup>1Topologically $`𝐂/𝐙_2𝐂`$. We keep the former to emphasize that this description also captures the asymptotic metric on the space. $`\|_{U(1)}𝐂/𝐙_2\times ^1`$ (3.23) Note that as $`z0`$, the $`^1`$ does not vanish. In this limit the equations (3.22) are solved by $`\|\omega \|^2=r`$ and $`a_i=0`$, while $`b_i`$ parameterize a $`^1`$ with Kähler class $`2r`$. Let us now turn to the submanifold describing vortices in orthogonal $`U(1)`$ subgroups as in (2.24). The vortices should now pass through each other without interacting. This submanifold is defined by the requirement $`\omega =0`$, and<sup>2</sup><sup>2</sup>2 This situation is similar to the theory describing two D-strings, in which $`r=0`$ and there is no $`\mathrm{\Psi }`$ field, forcing $`\omega =0`$ . we will refer to it as $`\|_{\omega =0}`$. The first two equations in (3.17) tell us that $`a_i`$ and $`b_i`$ each define a point on $`^1`$, while the third equation, which reads $`a_i\overline{b}_i=0`$, requires these points to be antipodal. Again, acting with the $`SU(2)_F`$ symmetry then sweeps<sup>3</sup><sup>3</sup>3Generically the orbits of the $`SU(2)_F`$ action on $`\stackrel{~}{M}_{2,2}`$ are three dimensional. They degenerate to two dimensional orbits on $`\|_{U(1)}`$ and $`\|_{\omega =0}`$. out a $`^1`$. We still have to divide out by the $`𝐙_2`$ gauge action which this time acts as in (3.20), exchanging $`zz`$ and, at the same time, mapping antipodal points on $`^1`$. We therefore have, $`\|_{\omega =0}{\displaystyle \frac{𝐂\times ^1}{𝐙_2}}`$ (3.24) But what happens at the origin? If we set $`\omega =z=0`$ then $`\mathrm{\Psi }`$ feels the full force of the restored $`U(2)`$ gauge symmetry, resulting in a unique solution to the constraints (3.1). This means that when the vortices live in orthogonal $`U(1)`$’s, as in (2.24), the internal space collapses as they approach each other! $`\|_{\omega =0}`$ can be thought of as the cone over $`(𝐒^1\times ^1)/𝐙_2`$ (which can alternately be described as the non-trivial $`𝐒^2`$ bundle over $`𝐒^1`$, or as the connect sum $`^3\mathrm{\#}^3`$). Note that the submanifold $`\|_{\omega =0}`$ is singular at $`z=0`$. This is an artifact of restricting attention to this subspace; the full manifold $`\stackrel{~}{}_{2,2}`$ should be smooth at the point $`\omega =z=0`$. How can we understand the result that the internal space collapses at the origin of $`\|_{\omega =0}`$ from the perspective of the soliton solutions? In fact, it is rather simple to see. The solutions of the form (2.24) generically transform non-trivially under the $`SU(2)_{\mathrm{diag}}`$ vacuum symmetry, sweeping out the $`^1`$ internal space. However, as $`x_1x_2`$ (corresponding to $`z0`$), the relevant part of $`A_\mu `$ and $`q`$ approaches the unit matrix, and the $`SU(N)_{\mathrm{diag}}`$ symmetry no longer acts. Two coincident vortices in orthogonal $`U(1)`$ sectors have no internal space! This is one of the important points that allows for reconnection to generically occur in this model. #### Reconnection and the Origin of Moduli Space Having determined the topology of these two submanifolds, let us now examine whether reconnection takes place on each. We start with $`\|_{U(1)}`$. Here the argument proceeds as for the abelian case: the submanifold $`𝐂/𝐙_2`$ is a smooth cone, as shown in Figure 2, with $`z^2`$ the single valued coordinate at the tip of the cone. This ensures that any trajectory hitting the tip of the cone results in right angle scattering. The $`𝐙_2`$ does not act on the internal space $`^1`$ and it plays no role in the discussion of reconnection. What about the space $`\|_{\omega =0}`$, describing orthogonal vortices? Here the issue is somewhat clouded by the singularity at the center of the space. However, consider first the resolved space where the $`^1`$ does not degenerate at the origin. Since the $`𝐙_2`$ action has no fixed points, such a manifold is smooth. In contrast to the previous case, a trajectory through the origin at $`z=0`$ now corresponds to vortices passing straight through each other; there is no right angle scattering. The reason for this is that the $`𝐙_2`$ gauge symmetry does not act only on $`𝐂`$ but also on the internal space; it exchanges both the positions and the identities of the particles. This means that, near the origin, $`z`$ is the single valued coordinate rather than $`z^2`$ and right-angle scattering does not occur. The true motion in $`\|_{\omega =0}`$, in which the $`^1`$ degenerates, can be thought of as the limiting case of this discussion. The need to take this degenerative limit is necessary since, as we shall presently see, $`\omega =0`$ corresponds to the only case of no reconnection; if the moduli space $`\|_{\omega =0}`$ were not described by this singular limit then, by continuity, vortices in the neighborhood of $`\omega =0`$ should also pass through each other. Let us now show that, as promised, vortices of arbitrary orientation always scatter at right angles unless $`\omega =0`$. As we have seen above, the key to this lies in the $`𝐙_2`$ action (3.16). In particular, we are interested in this action as the vortices approach each other and $`z0`$. Then we see that, provided $`\zeta =2z/\omega 0`$, the $`𝐙_2`$ action on $`\mathrm{\Psi }`$ and $`\omega `$ can be absorbed in the $`U(1)^2`$ gauge transformations, leaving only $`zz`$. In other words, for all $`\omega 0`$, the single valued coordinate near the origin is $`z^2`$ rather than $`z`$. Using the general arguments described above, this implies right-angle scattering and reconnection of cosmic strings. To complete the argument, we need to make sure that restricting to $`z=0`$ where the vortices coincide defines a complete submanifold. Let us denote this as $`\|_{z=0}`$. It may be defined in a coordinate independent manner as the fixed locus of the $`U(1)_R`$ action (since the resulting action on $`\omega `$ is gauge equivalent to a $`SU(2)_F`$ rotation). When $`z=0`$, the $`𝐙_2`$ action (3.16) can be absorbed into the axial combination of the $`U(1)_1\times U(1)_2`$ gauge symmetry (3.13) and $`\|_{z=0}`$ can be thought of as the resolution of the $`𝐙_2`$ fixed point. What is $`\|_{z=0}`$? Upon setting $`z=0`$, the constraints (3.17) read $`\|a_1\|^2+\|a_2\|^2+\|\omega \|^2=r,\|b_1\|^2+\|b_2\|^2\|\omega \|^2=r,a_i\overline{b}_i=0`$ (3.25) where we must still quotient by the $`U(1)^2`$ action (3.13). The first constraint, together with the $`U(1)_1`$ action, defines a copy of $`^2`$ (although it doesn’t inherit the round Fubini-Study metric). For a generic point $`p^2`$, the second and third constraints in (3.25) determine $`b_i`$ uniquely up to a phase, which is gauged away by $`U(1)_2`$. If this were true globally, we would have $`\|_{z=0}^2`$ (3.26) However, there are two exceptional points. Firstly, when $`a_i=0`$ and $`\|\omega \|^2=r`$, the $`b_i`$’s parameterize a $`^1`$ rather than a point. We have $`\|_{U(1)}\|_{z=0}^1`$. Secondly, when $`\omega =0`$, the full $`U(2)`$ gauge symmetry is restored and the $`a_i`$’s parameterize a point rather than a $`^1`$; we have $`\|_{\omega =0}\|_{z=0}\{0\}`$. In fact, it is a rather cute fact that after making these two adjustments, (3.26) remains correct! To see this, we may think of $`\|_{z=0}`$ as a fiber over the interval $`\|\omega \|[0,\sqrt{r}]`$. The phase of $`\omega `$ shrinks to zero at each end due to the action of the gauge symmetry. In the middle of the interval, $`a_i`$ and $`b_i`$, modulo the constraints (3.25) and the $`U(1)^2`$ gauge action, define a $`\mathrm{𝐂𝐏}^1`$. The phase of $`\omega `$ is fibered over this to yield a $`𝐒^3`$ (to see this, note that rotating the phase of $`\omega `$ is gauge equivalent to an $`SU(2)_F`$ flavor transformation). We therefore have a description of $`\|_{z=0}`$ in terms of an $`𝐒^3`$ fibration over the interval, degenerating to a point at one end and to $`^1`$ at the other. This is precisely $`^2`$. In summary, as two vortices collide their relative orientations define a point $`p^2`$. The vortices undergo $`90^\mathrm{o}`$ scattering (and, hence, strings undergo reconnection) unless $`p`$ coincides with the special point $`\omega =0`$ on $`^2`$. Thus the reconnection probability $`P`$ is unity. Note that we did not assume geodesic motion on the moduli space, and the argument for reconnection goes through even for strings carrying different currents in the internal space. Although this may seem surprising, similar results were observed numerically for Witten’s superconducting strings, resulting in excess charge build up at the interaction point . The behavior of the vortices passing through the special point $`\omega =0`$ presumably depends on the quantity $`\zeta =2z/\omega `$ as $`z0`$. For $`\zeta \mathrm{}`$, we have seen that the vortices pass through each other unscathed. We suspect that for other values of $`\zeta `$, the vortices undergo scattering an angle less than $`90^\mathrm{o}`$. ### 3.4 Reconnection of $`U(N)`$ Strings We now turn to the case of vortices in $`U(N)`$ gauge theories. The details are similar to the $`U(2)`$ case so we shall be brief. The matrix $`Z`$ remains $`2\times 2`$, while $`\mathrm{\Psi }`$ is now a $`2\times N`$ matrix. Once again we may employ the auxiliary $`U(2)`$ gauge action to place $`Z`$ in upper triangular form: $`Z=\left(\begin{array}{cc}z& \omega \\ 0& z\end{array}\right),\mathrm{\Psi }=\left(\begin{array}{ccc}a_1& \mathrm{}& a_N\\ b_1& \mathrm{}& b_N\end{array}\right)`$ (3.31) a choice which is preserved by the remnant $`U(1)_1\times U(1)_2\times 𝐙_2`$ action of equations (3.13) and (3.16). The same analysis of the previous section shows that asymptotically, $`\stackrel{~}{}_{2,N}{\displaystyle \frac{𝐂\times ^{N1}\times ^{N1}}{𝐙_2}}`$ (3.32) As we have seen, the question of reconnection boils down to the $`𝐙_2`$ action which, since it is unchanged, ensures that $`z^2`$ is the single-valued coordinate as $`z0`$ provided $`\zeta 0`$. This time the manifold $`\|_{z=0}`$ has complex dimension $`2N2`$, and is defined by the quotient construction, $`{\displaystyle \underset{i=1}{\overset{N}{}}}\|a_i\|^2+\|\omega \|^2=r,{\displaystyle \underset{i=1}{\overset{N}{}}}\|b_i\|^2\|\omega \|^2=r,{\displaystyle \underset{i=1}{\overset{N}{}}}a_i\overline{b}_i=0`$ (3.33) One can view this space as a fibration over the interval<sup>1</sup><sup>1</sup>1We’re grateful to James Sparks for discussions and explanations regarding these issues. $`\|\omega \|^2[0,r]`$. To ensure that the space $`\|_{z=0}`$ is smooth, one must check that spheres degenerate at either end of the interval, rather than a more complicated space. For $`\|\omega \|^20,r`$, we may use the $`U(1)_2`$ action to set the phase of $`\omega `$ to a constant. Then the constraints (3.33) define the Stiefel manifold $`V(2,N)U(N)/U(N2)`$ of orthonormal two-frames in $`𝐂^N`$. Dividing by the remaining $`U(1)_1`$ action, the fiber over a generic point is $`V(2,N)/U(1)_1`$. At the two end points of the interval some submanifold of this fiber degenerates. When $`\|\omega \|^2=r`$, so $`a_i=0`$, the constraints (3.33), together with the $`U(1)_2`$ action, ensure that the fiber shrinks to $`^{N1}U(N)/U(N1)\times U(1)`$. This means that the degenerating cycle at $`\|\omega \|^2=r`$ is $`[U(N2)\times U(1)]/[U(N1)\times U(1)]𝐒^{2N3}`$ (3.34) Meanwhile, at the other end of the interval, when $`\omega =0`$, the full $`U(2)`$ gauge action is restored. When combined with the constraints (3.33), this gives rise to the Grassmanian of complex two-planes in $`𝐂^N`$, which can be described as $`G(2,N)U(N)/U(N2)\times U(2)`$. At this end the degenerating cycle is $`[U(N2)\times U(2)]/[U(N2)\times U(1)]𝐒^3`$ (3.35) We therefore find that the cycles that degenerate at either end of the interval are indeed spheres, and the space $`\|_{z=0}`$ is smooth. In summary, the submanifold $`\|_{z=0}`$ is smooth and reconnection occurs unless the vortices collide over the complex codimension 2 submanifold $`\|_{z=0}\|_{\omega =0}G(2,N)`$, therefore $`P=1`$. ## 4 Symmetry Breaking and Quantum Effects So far we have discussed the situation in which the $`SU(N_f)`$ flavor symmetry of the model is unbroken. Since global symmetries are unlikely to be exact in Nature, in this section we discuss various mechanisms by which the flavor symmetry is broken and/or the internal modes on the string are lifted. ### 4.1 Symmetry Breaking and Monopole Pair Creation We start by introducing explicit symmetry breaking terms into the Lagrangian. We will present two examples in which the moduli space of vortices gets lifted, leaving behind $`N`$ different cosmic strings, each carrying magnetic flux in a different part of the gauge group. The simplest symmetry breaking term is a mass for the scalar fields $`q_i`$. After a suitable unitary transformation, we have the potential, $`V={\displaystyle \frac{\lambda e^2}{2}}\mathrm{Tr}\left({\displaystyle \underset{i}{}}q_iq_i^{}v^2\right)^2+{\displaystyle \underset{i}{}}m_i^2q_i^{}q_i`$ (4.1) For $`m_i^2\lambda e^2v^2`$, the theory still lies in the Higgs phase, with vacuum expectation value $`q_i^a=\sqrt{v^2{\displaystyle \frac{m_i^2}{\lambda e^2}}}\delta _i^a\mu _i\delta _i^a\text{(no sum over }i\text{)}`$ (4.2) The symmetry breaking pattern now becomes, $`U(N_c)\times SU(N_f)SU(N)_{\mathrm{diag}}U(1)_{\mathrm{diag}}^{N1}`$ (4.3) where the first, spontaneous, breaking occurs at the scale $`e^2v^2`$, while the second, explicit breaking, occurs at the scales $`m_i`$. With the $`SU(N)_{\mathrm{diag}}`$ broken, we can no longer sweep out a moduli space of vortex solutions as in (2.14) and the internal $`^{N1}`$ space is lifted. What remains are $`N_c`$ distinct vortex solutions in which the non-abelian field strength has non-vanishing component within only one of the diagonal $`U(1)U(N_c)`$. Let $`B`$ denote the adjoint valued magnetic field in the direction of the strings. Then we may embed an abelian vortex in the $`i^{\mathrm{th}}`$ $`U(1)`$ subgroup of $`U(N_c)`$ with $`B\mathrm{diag}(0,\mathrm{},1,\mathrm{},0)`$ (4.4) Such a vortex has tension $`T_i\mu _i^2`$. Since each of these vortices is supported by the same topological quantum number, $`\mathrm{Tr}B`$, only one of these strings is globally stable; the others may all decay into the string with the lowest tension. We shall discuss one such mechanism for this decay shortly. It is simple enough to alter our model to arrange for all $`N_c`$ strings to have the same tension. We introduce a new, adjoint valued, scalar field $`\varphi `$, with canonical kinetic term, and consider the potential, $`V={\displaystyle \frac{\lambda e^2}{2}}\mathrm{Tr}\left({\displaystyle \underset{i}{}}q_iq_i^{}v^2\right)^2+{\displaystyle \underset{i}{}}q_i^{}(\varphi m_i)^2q_i`$ (4.5) This is a variant on the potentials that appear in $`𝒩=2`$ SQCD. Unlike the potential (4.1), symmetry breaking in the pattern (4.3) now occurs regardless of the relative values of $`m_i`$ and (non-zero) $`v^2`$. The unique, gapped, vacuum is given by $`q_i^a=v^2\delta _i^a,\varphi =\mathrm{diag}(m_1,\mathrm{},m_N)`$ (4.6) In this theory we again have $`N_c`$ different vortices, with magnetic flux (4.4), but now with equal tension (2.4). (Note that if we also include an explicit mass $`M`$ for the adjoint scalar $`\varphi `$ then symmetry breaking only occurs for suitably small $`M`$, but the tensions of the vortices remain equal). Both deformations (4.1) and (4.5) result in $`N_c`$ different types of vortices, each embedded in a different, orthogonal, $`U(1)`$ subgroup of $`U(N_c)`$. This ensures that two strings colliding with energies $`E\mathrm{\Delta }m_i`$ fall into one of the two categories discussed in Section 2: either the strings are of the same type (i.e. inhabit the same $`U(1)`$) and they reconnect; or they are of different types, and pass through each other. Unlike the situation where the strings enjoyed an internal moduli space, there is no need for fine tuning to make the strings miss each other: the potential does the job for us. Therefore at energies smaller than the mass splittings $`\mathrm{\Delta }m_i`$, the classical probability for reconnection is $`1/N`$. At energies much larger than this, the masses are negligible and the probability increases to unity (at least whenever the moduli space approximation of the previous section is valid). ### Confined Monopoles In fact, quantum effects give rise to a finite probability for reconnection even for distinct strings. This can occur if one string turns into another through the creation of a confined monopole. Here we give an estimate of the magnitude of this effect. The presence of confined monopoles, acting like beads on the cosmic string, may have other interesting cosmological consequences as explored in . Strings living in different $`U(1)`$ subgroups are supported by the same topological invariant $`\mathrm{Tr}B`$, suggesting that they may transmute into each other. The change of the string, from one type to another, occurs by a kink on the string worldsheet which, from the four-dimensional perspective, has the interpretation of a confined magnetic monopole. These monopoles were described in and further explored in . Similar objects were previously discovered in $`𝐙_2`$ strings in . The mass of the kink on the worldsheet is $`M_{\mathrm{kink}}r\mathrm{\Delta }m_i{\displaystyle \frac{2\pi \varphi }{e^2}}M_{\mathrm{monopole}}`$ (4.7) which has the same parametric dependence as the mass of the unconfined magnetic monopole. (In the supersymmetric context the equality $`M_{\mathrm{kink}}=M_{\mathrm{monopole}}`$ is exact). Reconnection for different abelian strings requires the quantum pair creation of a monopole-anti-monopole on the string as shown in Figure 3. One can estimate the probability for reconnection to occur by treating the worldsheet dynamics as a $`d=1+1`$ quantum field theory. For simplicity let us model the reconnection of two almost static strings at incident angle $`\theta `$ by the shape shown in Figure 3. Then reconnection reduces the energy of the configuration by $`\mathrm{\Delta }V=4Ta\mathrm{tan}^2(\theta /2)+2M_{\mathrm{monopole}}`$ (4.8) The reconnected region is specified to be $`a<x<a`$ where $`x`$ is the worldsheet spatial coordinate. This is the same potential arising in electron-positron pair creation in a constant electric field in $`d=1+1`$, for which the electric field times the positron charge is now given by $`2T\mathrm{tan}^2(\theta /2)`$. The famous result by Schwinger , evaluating the bounce action of a circular loop in Euclidean space, gives the decay width as $`\mathrm{\Gamma }\mathrm{exp}\left({\displaystyle \frac{\pi M_{\mathrm{monopole}}^2}{2T\mathrm{tan}^2(\theta /2)}}\right)\mathrm{exp}\left({\displaystyle \frac{\pi ^2(\mathrm{\Delta }m)^2}{e^4v^2\mathrm{tan}^2(\theta /2)}}\right)`$ (4.9) This computation ignores the relative velocity of the strings, and is valid only for almost parallel strings for which the exponent is large (and negative). It would be interesting to better quantify the role of these confined monopoles for other impact parameters. ### 4.2 Quantum Effects Until now, much of our discussion has been purely classical. Indeed, we have chosen the four-dimensional symmetry breaking scale $`e^2v^2`$ to be suitably high so the theory is weakly coupled. Nevertheless, the theory on the vortex string is necessarily strongly coupled at low-energies: it is the two-dimensional $`^{N1}`$ sigma model. For now let us set the masses $`m_i`$ discussed in the previous section to zero, ensuring that the $`SU(N)_{\mathrm{diag}}`$ symmetry is exact. The resulting low-energy quantum dynamics on $`^{N1}`$ is well understood. The Mermin-Wagner-Coleman theorem guarantees that the ground state wavefunction spreads over $`^{N1}`$, resulting in a unique vacuum state for the string. More quantitatively , the size of the vortex moduli space evolves under RG flow, resulting in dynamical transmutation and a mass gap for the internal modes on the vortex string. The one-loop beta function leads to the strong coupling scale $`\mathrm{\Lambda }_{^{N1}}=\mu \mathrm{exp}\left({\displaystyle \frac{2\pi r(\mu )}{N_c}}\right)`$ (4.10) where, from equation (2.15), we have $`r2\pi /e^2`$ at the symmetry breaking scale $`\mu =ev`$. Thus the currents discussed previously, which classically may travel along the worldsheet at the speed of light, become massive and do not persist. One can show using the large N expansion that the theory confines and all dynamical degrees of freedom are singlets of $`SU(N)_{\mathrm{diag}}`$ . In terms of reconnection, the quantum effects do little to change the story: at low energies $`E\mathrm{\Lambda }_{^{N1}}`$, the strings lie in a unique ground state and reconnect with unit probability. At higher energies, $`\mathrm{\Lambda }_{^{N1}}<E<ev`$, asymptotic freedom of the sigma model ensures that the classical analysis of the previous section is valid and, once again, the strings reconnect. At energies $`Eev`$, numerical simulations are required to determine the issue. Introducing masses $`m_i`$ as in (4.1) or (4.5) leads to a weakly coupled theory when $`\mathrm{\Delta }m_i\mathrm{\Lambda }_{^{N1}}`$ and the results of the previous subsection hold only in this regime. ### 4.3 Fermionic Zero Modes The low-energy quantum dynamics of the string can be dramatically changed by the inclusion of fermionic zero modes . We may add Weyl fermions $`\xi `$ and $`\psi `$ to the bulk theory with Yukawa couplings of the schematic form, $`L_{\mathrm{Yukawa}}=\overline{\psi }\xi q`$ (4.11) Any such coupling will lead to chiral fermionic zero modes $`\chi `$ propagating on the vortex string. The exact nature of these zero mode depends on the properties of $`\xi `$ and $`\psi `$ and, as we now discuss, different color and flavor representations for $`\xi `$ and $`\chi `$ will lead to very different low-energy physics for the cosmic strings. Here we sketch two examples. More details will be given in a future publication. First consider the example in which $`\psi _i`$ transforms, like $`q_i`$, in the fundamental $`𝐍_𝐜`$ of the gauge group, as well as the fundamental $`𝐍_𝐟`$ of the flavor group (recall that $`N_c=N_f=N`$), while $`\xi `$ is a singlet under both. Then the Yukawa coupling (4.11) can be shown to give rise to a single chiral zero mode on the worldsheet with kinetic term, $`L_{\mathrm{zeromode}}=i\overline{\chi }\text{/}\chi `$ (4.12) Such zero modes do not couple to $`^{N1}`$ modes on the string and do not affect the low-energy dynamics. (Note that the four-dimensional anomaly can be cancelled by the addition of further fermions transforming in the conjugate representation which may give rise to further fermionic zero modes but do not qualitatively change the low-energy string dynamics). A more interesting example comes if we consider $`\psi `$ to transform, once again, in the $`(𝐍_𝐜,𝐍_𝐟)`$ representation of $`U(N_c)\times SU(N_f)`$, while $`\xi `$ transforms in the adjoint representation of $`U(N_c)`$ (and is a singlet under $`SU(N_f)`$). In this case index theorems ensure the existence of $`N`$ zero modes $`\chi _i`$ on the worldsheet. However, crucially, they now couple to the strongly interacting $`^{N1}`$ sector of the theory. Let $`\pi _i`$, $`i=1,\mathrm{},N`$ define homogeneous coordinates on $`^{N1}`$, such that $`_{i=1}^N\|\pi _i\|^2=r`$, with $`^{N1}`$ obtained after identifying $`\pi _ie^{i\alpha }\pi _i`$. Then the fermionic zero modes on the worldsheet can be shown to couple to a bosonic $`U(1)`$ current, $`L_{\mathrm{current}}=i\overline{\chi }_i(\pi _j^{}\stackrel{}{\text{/}}\pi _j)\chi _i`$ (4.13) Once again, the gauge anomaly can be cancelled by the addition of conjugate fermions in four dimensions. In fact, such action guarantees that the $`d=1+1`$ theory on the string is non-chiral, cancelling a related sigma-model anomaly on the worldsheet . So what is the consequence of the interaction (4.13)? The crucial point, as explained in , is worldsheet chiral symmetry breaking. Classically, the $`U(1)`$ chiral symmetry acts as $`\chi _ie^{i\beta \gamma _5}\chi _i`$ while, quantum mechanically, only a $`𝐙_{2N}`$ is non-anomalous. The strong coupling dynamics on the worldsheet induces a condensate for the zero modes, $`\chi \chi \mathrm{\Lambda }_{^{N1}}`$, breaking the discrete chiral symmetry yet further: $`𝐙_{2N}𝐙_2`$. We find a situation similar to that of Section 4.1, in which the moduli space of vacua is lifted, now at a scale $`\mathrm{\Lambda }_{^{N1}}`$, resulting $`N`$ different ground states. Recent studies of the low-energy dynamics of these vortex strings in both supersymmetric theories identify these ground states with the $`N_c`$ vortex strings lying in orthogonal, diagonal $`U(1)U(N_c)`$ subgroups . Once again we have a situation in which the strings reconnect with probability $`P=1/N`$ at energy scales $`E\mathrm{\Lambda }_{^{N1}}`$, and with unit probability at higher energies where the sigma-model becomes asymptotically free and the classical moduli space approximation holds. ## 5 Summary and Conclusions We have studied the low-energy dynamics of cosmic strings embedded in a $`U(N_c)`$ gauge theory with $`N_f=N_cN`$ scalars transforming in the fundamental representation. The cosmic strings in this theory obtain an internal $`^{N1}`$ space in which they move. We presented a number of deformations which lift this internal space at a scale $`M`$, leaving behind $`N_c`$ types of vortex string, each embedded in a different diagonal $`U(1)U(N_c)`$. Strings of the same type reconnect, while strings of different types do not interact. In this manner, the classical probability of two cosmic strings reconnecting at energies $`EM`$ is $`P=1/N`$. A reconnection probability of $`P=1/N`$ may also be achieved by simply considering $`N`$ decoupled abelian-Higgs models. Our strings are distinguished from this trivial case by two effects. Firstly, even at energies $`EM`$, the quantum creation of confined magnetic monopoles may lift the probability above $`P=1/N`$. We have only been able to compute this effect in the limit of small velocity and small angle of incidence where it is negligible. Away from this regime it may be the dominant contribution to reconnection and a better understanding of this process is desired. The second effect occurs at energy scales $`EM`$ where the classical probability for reconnection increases to $`P=1`$, at least when the moduli space approximation is valid. We find it interesting that semi-classical magnetic strings in non-abelian gauge theories remain strongly coupled at large $`N_c`$, distinguishing them from their non-perturbative electric counterparts in QCD-like theories, which are expected to interact with coupling $`1/N_c^2`$. For energies beyond the moduli space approximation we have been unable to determine the probability of reconnection analytically, although experience with the abelian-Higgs model suggests that it may remain unity up to very high energies. It would, of course, be interesting to develop numerical simulations to extract the functional dependence of the probability over the full ranges of impact velocities and incidence angles. Finally, an important, outstanding problem is to determine how the scaling of the string network is affected by the presence of a threshold scale $`M`$, and the associated confined monopoles, acting like beads on string, which are created after reconnection. One would expect monopoles of to also be created by the Kibble mechanism during formation of the initial string network. It seems plausible that a suitably chosen $`M`$ may skew the velocity distribution of the strings, giving rise to a larger concentration of low-energy strings. This would distinguish our non-abelian cosmic strings from others such as abelian strings ($`P=1`$), strongly coupled QCD-like strings ($`P1/N_c^2`$), weakly coupled fundamental strings ($`Pg_s^2`$) and D-strings/wrapped D-branes ($`P<1`$). One can only hope that cosmic strings are one day observed, presenting us with the challenge of deciding between these different possibilities. ## Acknowledgement Our thanks to Nima Arkani-Hamed, Nick Dorey, Amihay Hanany, Sean Hartnoll and James Sparks for many useful discussions. K. H. is supported in part by the Grant-in-Aid for Scientific Research (#12440060 and #15540256) from the Japan Ministry of Education, Science and Culture, and by funds provided by the U.S. Department of Energy (D.O.E.) under cooperative research agreement DF-FC02-94ER40818. D. T. is grateful to the Royal Society for funding.
warning/0506/astro-ph0506484.html	ar5iv	text	# Distant Field BHB Stars III: Identification of a probable outer halo stream at Galactocentric distance r=70kpc ## 1 Introduction The existence of a dark massive halo appears to be a generic feature of many galaxies, but the total masses, sizes, and the formation history of galactic halos are poorly understood. This is mostly because we do not have large enough samples of dynamical tracers at sufficiently large radii. The formation and extent of such mass distributions are of great importance in understanding the nature of the dark matter and its role in galaxy formation and evolution. For instance, the quantification of the dark matter content of the Galaxy would allow us to construct a picture of the assembly of the various baryonic components, through comparison with simulations. The baryonic components can, in turn, provide information about the evolution of the halo. There is compelling evidence that at least part of the stellar halo has been built up via the accretion of smaller satellite galaxies. Numerous searches have been made for streams of material responsible for building up the Galaxy. A striking example is the identification of the Sagittarius dwarf galaxy (Ibata, Gilmore & Irwin, 1994) and its stellar stream (e.g. Helmi & White, 2001). Recently, an extensive stream of stars has been uncovered in the halo of the Andromeda galaxy (M31), revealing that it too is cannibalising a small companion (e.g. Lewis et al. 2004). Such streams yield crucial information on the accretion history and formation of galaxy halos. Extended stellar streams have also been used to constrain the mass of the Galactic halo (e.g. Johnston et al. 1999) and in M31 (Ibata et al. 2004). A number of authors have noted the possible evidence for streams in the distribution of the intrinsically rare, but very luminous, carbon stars. Sanduleak (1980) proposed the association of a single carbon star with the Magellanic Stream and Totten & Irwin (2000) made the general observation that the non-uniform distribution of carbon-stars in their extensive survey of the halo may indicate that a number of the stars are associated with streams. We have previously argued (Clewley et al., 2002, hereafter Paper I) that blue horizontal branch (BHB) stars are an ideal population for exploring the outer reaches of the Halo. Like carbon stars they are luminous standard candles but are also far more numerous. BHB stars are A–type giants. A–type stars in the Galactic halo are easily identified, as they lie blueward of the main–sequence turnoff (e.g. Yanny et al., 2000, hereafter Y2000). Unfortunately assembling clean samples of remote r $`>60`$kpc <sup>1</sup><sup>1</sup>1In this paper we use the coordinate $`r`$ to denote Galactocentric distances and the coordinate $`R`$ to denote heliocentric distances. BHB stars is made difficult by the existence of a contaminating population of high–surface–gravity A–type stars, the blue stragglers, that are between one and three mag. fainter. Previous analyses required high signal-to-noise ratio ($`S/N`$) spectroscopy to reliably separate these populations (e.g. Kinman, Suntzeff, and Kraft, 1994), making identification of BHB stars in the distant halo unfeasible. This paper is the third in a series. In Paper I we developed two classification methods that enabled us to overcome the difficulties in cleanly separating BHB stars from blue stragglers, and outlined an observational programme to survey the halo for BHB stars. In the second paper (Clewley et al., 2004, hereafter Paper II), we presented photometry and spectroscopy of faint $`16.0<B<19.5`$ candidate BHB stars in two northern high Galactic latitude fields and four southern fields. This work resulted in a sample of 60 BHB stars at distances $`11<R<52`$kpc (mean $`28`$kpc), with measured radial velocities. Here we apply the methods of Papers I and II to survey for halo BHB stars at much greater distances, $`65<R<115`$kpc. The new survey uses Sloan Digital Sky Survey (SDSS) photometry to isolate a sample of faint halo A–type stars. Reliable classifications are derived from medium resolution spectroscopy using FORS1 at the VLT. The candidate BHB stars observed at the VLT were selected using $`u^{}g^{}r^{}`$ photometry from the SDSS Early Data Release (EDR) data set (Stoughton et al., 2002). The EDR photometry was preliminary, and has since been revised. Because we need accurate Johnson-Kron-Cousins $`B,V`$ magnitudes for the classification, we have used the more recent SDSS Data Release 2 (DR2) photometry (Abazajian, 2004) of the EDR–selected candidates for this purpose. In Section 2 of this paper we describe the selection of the BHB candidates from the EDR data set, and provide our prescription for transforming the DR2 $`g,r`$ magnitudes of A–type stars to $`B,V`$ magnitudes. Section 3 provides a summary of the VLT spectroscopic observations, and the data reduction and line measurement procedures followed. In Section 4 we use the methods of Papers I and II to classify these stars into categories BHB and blue straggler, and provide a summary table of distances and radial velocities of the eight stars classified as BHB. We compute the velocity dispersions of the two populations and compare them with previous work. In Section 5 we discuss the kinematics of the BHB stars. We perform an orbital analysis of the sample and suggest that most of them are plausibly associated with a single orbit. ## 2 Selecting the BHB candidates ### 2.1 Colour selection We selected candidate BHB stars using the SDSS EDR point spread function (PSF) $`u^{}g^{}r^{}`$ magnitudes of stellar objects. The SDSS photometry has evolved between EDR and DR2 for a variety of reasons: i) the reference photometric system is now that of the SDSS 2.5m telescope itself, rather than the photometric monitoring telescope, ii) the absolute calibration of the standards has improved, and the calibration of the survey data relative to the standards has improved, iii) instrumental systematics (especially scattered light) are now better understood. Therefore, we have taken advantage of the improved photometry in DR2 for the subsequent analysis, in particular in classifying the objects and estimating their distances. To distinguish between photometric systems, the EDR system is designated by asterisks, the system of the photometric monitoring telescope is designated by primes, and the DR2 magnitudes are unadorned, which is the SDSS convention. Bearing in mind that all the stars are in the remote halo, all the SDSS magnitudes discussed in this paper have been corrected for Galactic extinction, using the map of Schlegel, Finkbeiner & Davis (1998). We limited ourselves to the northern equatorial stripe in the EDR data set, which covers $`145^{}<\alpha <236^{}`$, $`1.25^{}<\delta <+1.25^{}`$ (J2000). This is SDSS stripe 10, observed in runs 752 and 756 (Stoughton et al., 2002). In selecting candidate distant BHB stars from the EDR, we were guided by the results of Y2000, who studied the spatial distribution of a sample of A–type stars selected from the EDR using the (reddening–corrected) colour selection box $`0.3<g^{}r^{}<0.0`$, $`0.8<u^{}g^{}<1.5`$. In plotting apparent magnitude against $`\alpha `$ for these A–type stars, Y2000 discovered that the distribution is not smooth, but includes striking over–dense regions occurring in bands. These have subsequently been identified as tidal debris from the Sagittarius dwarf galaxy. Furthermore the bands occur in pairs, coincident on the sky, but separated by $`2`$mag. The brighter objects are the more luminous BHB stars, and the fainter objects are blue stragglers at the same spatial location in the halo. Our aim was to use the EDR data set to select a sample of candidate faint BHB stars, with minimal contamination by blue stragglers, in order to make the most efficient use of the spectroscopic time awarded (for classification and velocity measurement). To reach large distances, we chose the magnitude range $`20.0<g^{}<21.1`$, corresponding to $`65<\mathrm{R}<115`$kpc, if the objects are BHB stars. Fortunately, as demonstrated by Lenz et al. (1998) using synthetic photometry, the $`u^{}g^{}r^{}`$ colours of A–type stars show some dependence on surface gravity. Given the accuracy of the photometry at the distances of interest, $`0.1`$mag., the $`u^{}g^{}r^{}`$ colours cannot provide reliable separation of the two populations. Nevertheless Y2000 demonstrated that a colour cut in the $`u^{}g^{}`$ verus $`g^{}r^{}`$ plane is effective in enhancing the contrast of the individual bands of tidal debris i.e. can substantially reduce the contamination of one population by the other. Another way of looking at this is illustrated in Fig. 1. The upper diagram plots the colours of all EDR stars within the colour selection box of Y2000, for the limited range $`200^{}<\alpha <230^{}`$, which is the region of the EDR containing the strongest tidal debris bands. Stars in the brighter band $`18.8<g^{}<19.2`$, which should be predominantly BHB stars, are marked with solid symbols, while stars in the fainter band $`20.5<g^{}<21.5`$, which should be predominantly blue stragglers, are marked with open symbols<sup>2</sup><sup>2</sup>2The larger magnitude interval selected is because blue stragglers have a larger spread in luminosity than BHB stars (§3.2).. For reference the dashed line shows the dividing line used by Y2000. The brighter band, typical colour error $`\sigma (u^{}g^{})=0.05`$, is mostly confined to a narrow range in $`u^{}g^{}`$, which evidently defines the colour domain of the BHB stars. The fainter stars, open symbols, are concentrated towards the top of the plot, but are spread over a larger colour range. However, much of the spread is accounted for by the larger colour errors $`\sigma (u^{}g^{})=0.2`$, as can be seen by reference to the histogram in the lower plot. It is evident that the mean $`u^{}g^{}`$ colour of blue stragglers is substantially bluer than for BHB stars. The spectroscopic classification criteria (detailed in §4) work best near $`(BV)_0=0.1`$, which corresponds to $`g^{}r^{}=0.125`$, using the colour transformation provided by Fukugita et al. (1996). On this basis we adopted the colour cuts shown by the box in Fig. 1, defined by $`1.08<u^{}g^{}<1.40`$, $`0.2<g^{}r^{}<0.04`$, and limited candidate selection to objects with colour error $`\sigma (g^{}r^{})<0.07`$. While these colour cuts should be nearly optimal in terms of the fraction of candidates that are BHB stars, we would still expect substantial contamination by blue stragglers, on account of the large $`u^{}g^{}`$ colour errors at the faint magnitudes of the sample, $`20.0<g^{}<21.1`$. The distribution in $`\alpha `$ and $`g^{}`$ of all the stars satisfying these criteria is shown in Fig. 2. The higher density of points at $`\alpha >200^{}`$ is due to blue stragglers in the Sagittarius tidal stream. Our classification methods produce samples of BHB stars that are contaminated by blue stragglers at the level of $`<10\%`$, for samples of A–type stars with a typical mix of the two populations (Paper I). The contamination of our BHB sample would be substantially greater if we attempted to classify stars in this region, so we confined our sample to $`\alpha <200^{}`$. The clump visible at $`\alpha =153^{}`$, $`g^{}20.2`$, is the horizontal branch of the Sextans dwarf spheroidal <sup>3</sup><sup>3</sup>3Note that the discussion in Y2000 of these objects mistakenly cites a paper concerned with a different galaxy, Sextans A., centre $`\alpha =153.3^{}`$, $`\delta =1.61^{}`$, J2000 (Irwin and Hatzidimitriou, 1995). In order to avoid stars in Sextans we confined our selection to $`\alpha >160^{}`$. The final selection includes 54 objects, of which we observed the 35 listed in Table 1. One of these, no. 32, proved to be a quasar. Looking ahead, of the remaining 34 candidates we are able to classify 20. These classifications are indicated on Fig. 2, with solid circles representing the eight stars classified BHB, and open circles representing the 12 stars classified blue straggler. For the remaining 14 candidates the classifications are uncertain, mostly due to insufficient $`S/N`$. In Table 1 we provide details of the 35 objects observed. Column 1 is our running number, and column 2 lists the coordinates. Successive columns provide the dereddened DR2 $`g`$ magnitude, and the dereddened $`ug`$ and $`gr`$ colours. The last column provides the dereddened $`BV`$ colour, calculated using the transformation derived in §2.2. There is good agreement between the EDR and DR2 photometry in the mean, but with noticeable scatter. For example, looking at the difference $`(g^{}r^{})(gr)`$ for our targets, the mean is 0.00 mag., and the standard deviation is 0.03 mag. ### 2.2 Transforming from $`gr`$ to $`BV`$ We use two methods to classify stars into categories BHB star and blue straggler. One method makes use of $`(BV)_0`$ colours. Therefore we need to convert the extinction corrected $`gr`$ colours to $`BV`$. Fukugita et al (1996) quote the transformation $$BV=(g^{}r^{}+0.23)/1.05,$$ (1) computed from synthetic photometry. Smith et al. (2002) quote the transformation $$BV=(g^{}r^{}+0.19)/0.98,$$ (2) based on observations of standard stars. These two relations are plotted in Fig. 3, as the dashed and dotted lines respectively. The two linear relations differ somewhat even over the narrow range of colours of interest in this paper, $`0.2<gr<0.04`$, by some 0.03 mag. Whereas one would naturally prefer the empirical relation over the synthetic relation, there is some indication that the actual transformation is non–linear in the region of the A stars, as the four stars measured by Smith et al. (2002) in the colour range of interest, lie on average 0.04 mag. above the linear relation (which is a fit over a wide colour range). Further evidence that the relation is non–linear comes from our own photometry, which is less precise, but has many more stars. Plotted in Fig. 1 are the colours of the 60 stars in the colour range $`0.05<(BV)_0<0.40`$, with $`(BV)_0`$ measured by ourselves (Paper II), which also have SDSS DR2 $`gr`$ colours. Only the errors on $`(BV)_0`$ are plotted, as these dominate. Our data points are systematically high relative to the linear relation of Smith et al. (2002). An additional source of uncertainty is the fact that the above linear relations were derived for the $`u^{}g^{}r^{}i^{}z^{}`$ system of the photometric monitoring telescope, slightly different from the $`ugriz`$ system of the SDSS 2.5m telescope itself, used for DR2. Accurate $`(BV)_0`$ photometry is required for the classification of the stars using the $`D_{0.15}`$–Colour method (§4). A systematic error in the $`(BV)_0`$ colour as large as 0.05 mag. could result in many of the classifications being in error. Therefore we re-investigated the colour transformation. We computed synthetic colours using the methods detailed in Hewett et al. (in prep.). We used model stars with \[Fe/H\]$`=1`$ and log $`g`$=3.5 (an appropriate surface gravity midway between BHB stars and blue stragglers), from Kurucz (1993). We fit a cubic polynomial to the relation between the $`(BV)_0`$ and $`(gr)_0`$ synthetic colours, for the colour range of Fig. 3. Finally, bearing in mind the uncertainty in the absolute calibration of the SDSS magnitudes onto the AB system (Fukugita et al., 1996), we allowed the zero point of the relation to be a free parameter, established by shifting the derived curve vertically to give the best fit to the data of Fig. 3. The curve provides a better fit than the two linear relations plotted. Furthermore the average offset of the four stars measured by Smith et al. (2002) in the colour range of interest reduces to 0.005 mag. We have therefore adopted this relation. The transformation is given by $$BV=0.764(gr)0.170(gr)^2+0.715(gr)^3+0.218,$$ (3) and is plotted in Fig. 3 as the bold solid line. We stress that this relation is specifically for A–type stars, and is not expected to be reliable for other types of star. The adjustment to the zero point was very small, only 0.014 mag. The agreement is encouragingly good and gives considerable confidence in the many elements going into this comparison – the SDSS DR2 photometry, the absolute calibration of Vega, and of the SDSS standard stars, our own photometry, the measurement of the different passband response functions, and the synthetic stellar spectra. The computed $`(BV)_0`$ colours of our targets are listed in the final column of Table 1. In summary, we have presented evidence that, in the colour range of interest, A stars lie systematically off the simple linear colour transformation measured by Smith et al. (2002), by about 0.05 mag. and we have derived a cubic relation that provides an improved fit. ## 3 Spectroscopic observations, analysis, and results ### 3.1 Observations We used the VLT FORS1 instrument, in service mode, over two periods, from 2003/03/25 to 2003/04/09, and from 2004/01/30 to 2004/03/20, to obtain medium resolution optical spectra of the 35 BHB candidates. The instrument is equipped with a $`2048^2`$ Tek CCD, with a projected scale of 0.2″pixel<sup>-1</sup>. We used the 600B grating, giving a dispersion of 1.2 Å pixel<sup>-1</sup>. With the 0.7″slit, the resolution achieved, measured from arc lines, was about 4Å, which is sufficient for the line–fitting procedure. The spectral coverage was 3400–5700 Å, which includes the relevant lines H$`\delta `$, H$`\gamma `$, and Ca II K $`\lambda 3933`$Å. Three BHB radial velocity standard stars in the globular cluster M5 were observed twice each. Table 2 summarises relevant information on the standards. Columns (1) to (5) list the identification, RA and Dec., $`V`$ magnitude, $`(BV)_0`$ colour, and the heliocentric radial velocity, V. The information in successive columns (6) to (11) in Table 3 contains averages of $`H\delta `$ and $`H\gamma `$ line measurements. Columns (6) to (8) in Table 2 list, respectively, the parameters $`D_{0.15}`$, $`b`$, and $`c`$ (explained in §4). The errors on the parameters $`b`$ and $`c`$ are provided in columns (9) to (11) in the form of $`A`$ and $`B`$, the semi-major and semi-minor axes of the error ellipse in the $`bc`$ plane, and $`\theta `$ the orientation of the semi-major axis, measured anti-clockwise from the $`b`$-axis. Here the error corresponds to the $`68\%`$ confidence interval for each axis in isolation (see Paper I for further details). The requested integration times of between 765 and 1980 seconds, for stars in the magnitude range $`20.0<g^{}<21.1`$, were estimated using the FORS1 exposure-time calculator, on the basis of the requested seeing, transparency, and lunar phase, in order to achieve the minimum continuum $`S/N`$ ratio of 15 $`\mathrm{\AA }^1`$ required to classify the stars (Paper I). All targets were observed near culmination, with a mean airmass of $`1.20\pm 0.1`$. In the event, observations of 12 of the 34 targets failed to achieve the required $`S/N`$, and therefore these targets cannot be reliably classified. The failures were primarily in the cases where the seeing was poor. In retrospect, for these service observations we should have included a safety margin in the requested integration times. All image frames were automatically bias and flat–field corrected by the FORS pipeline, and we then followed standard procedures for sky subtraction, spectral extraction, and wavelength calibration. An error spectrum, used for the line profile fits, was computed from Poisson considerations. For the wavelength calibration, HgCdHe arc observations were made during the day, and were used to derive the dispersion solution. To account for any flexure of the instrument, the zero point of the dispersion solution was established for each spectrum using the \[OI\] night-sky line at 5577.34 Å. We found a rms drift of the zero point of $`13`$kms<sup>-1</sup> over the entire data set. ### 3.2 Analysis As stated earlier, one of the candidates proved to be a quasar. In the remainder of the paper we ignore this object, and refer only to the 34 stars. The spectra were used to measure the shapes and widths of the H$`\delta `$ and H$`\gamma `$ lines (for classification), the EW of the CaII K line (to measure the metallicity), and the radial velocity of each star (for future dynamical analysis). For these measurements we followed the procedures set out in Papers I and II exactly, and we refer the reader to those papers for full details. Below we provide a brief explanation of how the Balmer lines were measured, and how the CaII K line EW is used to determine the metallicity. We then summarise how the magnitude, colour, and metallicity are combined with the classification to estimate a distance for each star. Balmer line profiles. After normalising each spectrum to the continuum, we fit a Sérsic function, convolved with a Gaussian of FWHM the instrumental resolution, to the H$`\delta `$ and H$`\gamma `$ lines. Two parameters of the fit, the scale width $`b`$, and the shape index $`c`$, are recorded. One classification procedure, the Scale width–Shape method, plots these two quantities against each other. A third quantity $`D_{0.15}`$, which is the line width at a depth $`15\%`$ below the continuum, is derived from $`b`$ and $`c`$. The second classification method, the $`D_{0.15}`$colour method, plots this quantity against $`(BV)_0`$. Because $`D_{0.15}`$ is a function of $`b`$ and $`c`$, the two classification methods are not completely independent. Metallicities from CaII K lines. The CaII K line is the strongest metal line present in the wavelength range covered by the spectra, and the only useful line in moderate resolution blue spectra for measuring metallicity. Plotting CaII K line EW against $`(BV)_0`$, the metallicity is determined by interpolation between lines of constant metallicity on this diagram (see Fig. 4). The uncertainty is established from the uncertainties of the two quantities plotted, and an additional uncertainty of 0.3dex is added in quadrature. This is the systematic error, and was established by comparing metallicities derived by this method using high $`S/N`$ data of comparable resolution, with accurate metallicities determined from high–resolution spectra. No attempt has been made to remove the possible contribution of interstellar CaII K absorption from the stellar K measurements. For a remote halo star the typical CaII K EW is 0.11Å/sin$`b`$ (Bowen, 1991), with a $`95\%`$ range of $`(0.060.31)`$Å/sin$`b`$. This range translates to $`0.070.35`$Å for our fields. Our 34 stars have mean EW 1.7Å and standard deviation 0.8Å, with only two stars having EW below 0.7Å. Therefore interstellar CaII K absorption is insignificant for the majority of our targets, but could bias the measured metallicities high for the small fraction of stars with the weakest lines. Distances. The absolute magnitude of a BHB star $`M_V(BHB)`$ depends on both metallicity and colour (i.e. temperature). In Paper II we derived a relation for the absolute magnitude of BHB stars by combining published relations for the dependence of $`M_V`$ on metallicity and on $`(BV_0)`$ colour, with the measured absolute magnitude at fixed metallicity and colour, as follows. The slope of the relation between apparent magnitude and metallicity for RR Lyrae stars in the Large Magellanic Cloud was measured by Clementini (2003) from observations of some 100 stars. We combined this with the measurement of the absolute magnitude of RR Lyrae stars at fixed metallicity, determined by Gould & Popowski (1998) from Hipparcos statistical parallaxes, to derive the linear relation for RR Lyrae stars: $$M_V(RR)=1.112+0.214[\mathrm{Fe}/\mathrm{H}].$$ (4) We then adopted a cubic expression determined by Preston et al. (1991), for the $`(BV)_0`$ colour dependence of the difference in absolute magnitudes between BHB and RR Lyrae stars, to produce the final expression for the absolute magnitude of BHB stars: $`M_V(BHB)`$ $`=`$ $`1.552+0.214[\mathrm{Fe}/\mathrm{H}]4.423(BV)_0`$ (5) $`+17.74(BV)_0^235.73(BV)_0^3.`$ Distances and associated errors are then determined using the apparent magnitudes $`V_0`$, and the corresponding photometric and metallicity errors. To compute $`V`$, we used the relation $`V=g^{}0.53(g^{}r^{})`$ (Fukugita et al., 1996), here disregarding the subtle differences between the different SDSS magnitudes ($`g,g^{},g^{}`$, etc.). The result produces distance errors of $`610\%`$ for our confirmed BHB stars. The exact form of Equation 5, particularly the zero point, remains controversial. Currently there are at least ten methods of determining the absolute magnitudes of RR Lyrae stars. We refer the interested reader to a recent review of this subject by Cacciari & Clementini (2003).<sup>4</sup><sup>4</sup>4They find, by averaging over all ten methods in their review, $$M_V(RR)=0.93\pm 0.12+(0.23\pm 0.04)[\mathrm{Fe}/\mathrm{H}]$$ (6) The absolute magnitudes of blue stragglers have been less well studied. Since we will not use the blue stragglers in any dynamical analysis, their distances are less interesting. In Paper II we adopted the following relation derived by KSK from data for blue stragglers in globular clusters published by Sarajedini (1993) $$M_V(BS)=1.32+4.05(BV)_00.45[\mathrm{Fe}/\mathrm{H}].$$ (7) ### 3.3 Results The results of these measurements for the 34 candidate BHB stars are provided in Table 3. Column (1) gives the number of the star, and columns (2) and (3) record the EW of the H$`\gamma `$ line, and the spectrum continuum $`S/N`$ per Å. Our classification methods were developed specifically for objects with strong Balmer lines, defined by EW H$`\gamma >13`$Å, and with continuum $`S/N>15`$Å<sup>-1</sup>. In all, only 20 of the 34 candidates meet both criteria, and are therefore classifiable. The majority fail because of inadequate $`S/N`$. The information in successive columns (4) to (9) in Table 3 contain averages of $`H\delta `$ and $`H\gamma `$ line measurements. The quantities in columns (4) to (9) in Table 3, i.e. $`D_{0.15}`$, $`b`$, and $`c`$, $`A`$, $`B`$ and $`\theta `$, are also provided for the radial velocity standards in columns (6) to (11) of Table 2. Comparing these quantities enables us to use the observations of the standard stars as a further check of our classification methods. In column (10) of Table 3 is listed the EW of the CaII K line. This is plotted against $`(BV)_0`$ in Fig. 4. The measured metallicity for each star, and the error (including random and systematic contributions) is provided in column (11). The large errors are a consequence of the comparatively large errors in the $`g^{}r^{}`$ colours. The mean measured metallicity of the stars plotted is $`1.4`$ with dispersion 0.6, similar to the mean value measured for our sample of brighter A–type stars (Paper II). There are no significant outliers, but this is not a strong statement, given the large errors. The radial velocity, corrected to the heliocentric frame, is provided in column (12), and the estimated distance, based on the classification from the following section, is provided in column (13). ## 4 Classification and velocity dispersion ### 4.1 Classification As noted above, of the 34 candidates, only 20 meet the requirements on spectroscopic $`S/N`$ and H$`\gamma `$ EW for reliable classification. In the following we restrict our discussion to the classification of these 20 objects. Of the other 14 candidates, five stars have EW H$`\gamma <13`$Å, and are considered unclassifiable. For the remaining nine candidates the $`S/N`$ of the spectra is too low for the classification to be reliable. We have nevertheless followed the classification procedures for these objects, but for clarity have omitted them from Figures 4 and 5. The final classifications are flagged as questionable. We have followed the classification procedures of Paper II (which are slightly different from those of Paper I) exactly, with the exception that we weight the two classification methods unequally, as detailed below. In Figure 5 we plot the two diagnostic diagrams for the 20 classifiable stars in the survey. The two figures are explained as follows. Figure 5(a) shows the $`D_{0.15}`$colour method. The average values of $`D_{0.15}`$ for H$`\gamma `$ and H$`\delta `$ against $`(BV)_0`$ are plotted for the 20 candidates. In Paper I we showed that reliable classification by this method requires the uncertainty on $`(BV)_0`$ to be less than 0.03 mag. Unfortunately this is untrue for most of the stars in our sample (Table 1). For this reason we give this method lower weight in the final classification. Figure 5(b) shows the Scale width–Shape method. The line–profile quantities $`b`$ and $`c`$, averaged for H$`\gamma `$ and H$`\delta `$ are plotted. The solid lines show the classification boundaries, from Paper II, with high–surface gravity stars (i.e. main–sequence A stars or blue stragglers, hereafter A/BS) above the line, and low–surface gravity stars (i.e. BHB stars) below the line. In both plots stars classified BHB are plotted as solid symbols and stars classified A/BS are plotted open. As we discuss below, the three triangles are stars that have ambiguous classifications, i.e stars that are classified as BHB by one classification method and not the other. Inspection of Figure 5 provides the following information. Of the 20 candidates, eight are classified BHB by the $`D_{0.15}`$colour method. The Scale width–Shape method classifies seven stars as BHB. There are six stars classified BHB by both methods. A total of nine stars are classified BHB by one or other of the methods. There is clearly close agreement between the two classification methods, but there are three stars with ambiguous classifications. Before considering these further, we note that the three radial velocity standards (Table 2), previously classified BHB from high–resolution spectroscopy, are all unambiguously classified BHB in both plots. The uncertainties on each parameter define the 2D probability distribution functions for any point. By integrating these functions below the classification boundary we can compute a probability $`P(BHB)`$ that any star is BHB. We can then average the probabilities for the two classification methods, to improve the classification. We have computed these probabilities for each star, giving twice the weight to the Scale width–Shape method when averaging (because, as mentioned above, the $`D_{0.15}`$colour method is affected by the relatively large colour errors). As in previous papers, stars with $`\overline{P}(BHB)>0.5`$ are then classified BHB, and stars with $`\overline{P}(BHB)0.5`$ are classified A/BS. Based on the Monte Carlo simulations of Paper I, we would expect the sample of BHB stars defined in this way to be contaminated by A/BS stars at no more than the $`10\%`$ level, which we consider satisfactory. Col. (14) of Table 3 provides the averaged probabilities, and corresponding classifications, for the 20 classifiable stars. Of the three stars with ambiguous classifications, 2 are classified BHB. The third star classified A/BS, star 9, is the object with the smallest value of $`D_{0.15}`$, and the smallest value of $`c`$. The small value of $`c`$ indicates a colour substantially redder than (although compatible with) the measured value of $`(BV)_0=0.16\pm 0.04`$. We also provide the classification probabilities for the nine stars with inadequate spectroscopic $`S/N`$, but for these the classifications are given as BHB? or A/BS? to indicate that they are not reliable. Finally the five stars with EW H$`\gamma <13`$Å are labelled unclassifiable. ### 4.2 Velocity dispersion Table 4 contains a summary of the kinematic properties of the final sample of eight BHB stars. Listed there are the Galactic coordinates $`l`$ and $`b`$, and the Galactocentric radial velocity and distance, V<sub>gal</sub> and $`r`$ respectively. To convert the heliocentric quantities to Galactocentric quantities, the heliocentric radial velocities are first corrected for solar motion by assuming a solar peculiar velocity of ($`U,V,W`$) = (-9,12,7), where $`U`$ is directed outward from the Galactic Centre, $`V`$ is positive in the direction of Galactic rotation at the position of the Sun, and $`W`$ is positive toward the North Galactic Pole. We have assumed a circular speed of 220 km s<sup>-1</sup> at the Galactocentric radius of the Sun ($`R_{}=8.0`$kpc). Table 4, then, distills the main observational result of the paper, a sample of distant BHB stars with measured radial velocities. The Table also includes three carbon stars, designated by their coordinates, which are introduced in Section 5. We find, after quadratically subtracting the measurement errors in the same manner as Norris & Hawkins (1991), that the measured dispersion of the radial component of the Galactocentric velocity dispersion for our BHB sample is 58$`\pm `$15km s<sup>-1</sup> at a mean heliocentric distance of $`80`$kpc. In Table 5 we compare this value against the measured velocity dispersion of a variety of samples. The sample of remote BHB stars is referred to as Sample A, and listed in the first line of Table 5. The first comparison sample, Sample B, comprises the 60 BHB stars $`11<R<52`$kpc, mean $`R=28`$kpc, from Paper II, which is the largest sample at such distances. Sirko et al. (2004a) have also isolated large samples of distant BHB stars using the SDSS. They split their sample into a bright ($`g<18`$) subsample, which is contaminated by blue stragglers at the level of about 10% (i.e. similar to the work presented here), and a faint subsample ($`g>18`$), which is contaminated at about 25%. If we consider only their clean bright sample, here Sample C, then $`\sigma =99.4\pm 4.3`$km s<sup>-1</sup> (Sirko et al. 2004b), at mean distance $`16`$kpc. Sample D consists of the 12 stars in Table 3 classified BS, with measured velocity dispersion $`129\pm 26`$km s<sup>-1</sup>, at mean distance $`40`$kpc. Finally considering the Galactic satellites discussed in §1, selecting the nine satellites within the distance range of our remote BHB sample, i.e. $`65<\mathrm{R}<102`$kpc, Sample E, we measure a velocity dispersion $`134\pm 32`$km s<sup>-1</sup>, at a mean distance $`82`$kpc. The velocity dispersions, mean distances, and sample sizes, of these four samples are entered in columns $`24`$ of Table 5. In the final column we list the probability that the measured value for the remote BHB stars could be drawn from the same population as each of the four comparison samples, as measured by the F–test. At better than the $`95\%`$ significance level, our sample of remote BHB stars has smaller velocity dispersion than BHB stars at much smaller radii ($`16`$kpc, Sample C), BS stars of similar apparent magnitude at intermediate radii ($`40`$kpc, Sample D), and satellites at the same radii ($`82`$kpc, Sample E). Compared to the sample of BHB stars at $`28`$kpc, the difference is not significant. We conclude that the velocity dispersion of the remote BHB stars is anomalously low, and in the following section we seek an explanation. ## 5 Discussion In considering the anomalously low velocity dispersion of the remote BHB stars, we first check that the significantly different velocity dispersion between the BHB and BS stars is robust. In Fig. 6 we plot classification probability against $`V_{gal}`$. In this plot, the large symbols represent the 20 stars with spectra of high S/N. BHB stars have $`\overline{P}>0.5`$, and, as usual, are shown by filled symbols, with A/BS stars marked by open symbols. Also plotted, as small symbols, are the nine stars with unreliable classifications. Finally the five unclassifiable stars are marked by small crosses, at $`\overline{P}=0.5`$. This plot shows a clear difference in the kinematics of stars at the bottom of the diagram (high probability A/BS, large velocity spread), and at the top of the diagram (high probability BHB, small velocity spread). There is no evidence for any BHB stars with large values of $`\|V_{gal}\|`$ that have been missed, because they fall just outside the classification boundary, or because they were unclassifiable because the spectra are of low S/N. Therefore, the difference in velocity dispersion between the two populations is quite robust to the method of classification. Another concern we had was the possibility that the sample is contaminated by misclassified blue stragglers in the Sagittarius stream. Fig. 2 plots $`\alpha `$ against $`g^{}`$ of the initial list of candidate BHB stars. As we discussed earlier this colour–selected sample of candidates is expected to be contaminated by blue stragglers, because of the large photometric errors at these faint magnitudes, and this is apparently confirmed by the high–density of stars at $`\alpha >200^{}`$, where presumably most of the stars are blue stragglers. It was therefore worrying that the eight BHB stars, marked in the upper diagram by filled circles, mostly lie close to the boundary $`\alpha =200^{}`$ in this plot, whereas one might expect them to be more uniformly scattered over the RA range. If these stars are misclassified blue stragglers, on the edge of the Sagittarius stream, this would provide a natural explanation for the small velocity dispersion. However, if there are any Sagittarius blue stragglers $`\alpha <200^{}`$ in our candidate list, most will be classified blue straggler, and the reduction in velocity dispersion would be greatest in our A/BS sample – the opposite of what is seen. We conclude from the foregoing discussion that we have succeeded in defining samples of BHB and A/BS stars, with small contamination, that show distinct kinematic properties. Indeed the fact that the measured velocity dispersions of the two populations are significantly different is confirmation of the reliability of the classification methods. In seeking an explanation for this difference, a number of possible dynamical explanations could be pursued, for example that the stellar orbits change from being predominantly radial to predominantly circular at large radii (e.g. Sommer-Larsen et al. 1997). While this is possible, a more convincing explanation was immediately apparent. Six of the BHB stars are confined to a small region of space with $`190^{}<\alpha <200^{}`$, $`63<r<78`$kpc (and a small range in $`\delta `$). The average distance of these stars from the centre of the Galaxy is $`70.6`$kpc. These stars are therefore confined to a very small fraction of the volume surveyed. The velocity dispersion of these six stars (corrected for measurement errors), $`42\pm 12`$km s<sup>-1</sup>, is too large to associate them with a bound object i.e. a low surface–brightness dwarf galaxy. The key to understanding the anomalous velocity dispersion is provided in Figure 7(upper), which plots Galactocentric radial velocity $`V_{\mathrm{gal}}`$ versus RA for these six BHB stars, marked by squares. A correlation is evident, indicative of streaming motion, perhaps associated with a disrupted satellite. Figure 7(lower) plots the position on the sky of the six stars. To investigate further the possibility of a stream, we searched the catalogue of faint carbon stars of Totten & Irwin (1998) in the vicinity, confining ourselves to the ranges RA $`160^{}<\alpha <200^{}`$, Dec. $`\pm `$ $`5^{}`$, and to similar distances (improved distance estimates taken from Totten, Irwin & Whitelock, 2000). Three stars meet these criteria, and have been added to both diagrams in Figure 7, marked as triangles. Remarkably the three stars appear to add to the evidence of a stream. Details of the three stars are provided in Table 4. (The velocity errors, taken from Totten, Irwin & Whitelock (2000), are 4, 4 and 6 km s<sup>-1</sup> in the order the stars appear in the Table) The correlation between $`V_{\mathrm{gal}}`$ and RA evident in Figure 7(upper), encompasses $`V_{\mathrm{gal}}=0`$, which would correspond to a turning point in the orbit at RA of $`195^{}`$. In order to investigate this trend in more detail, we consider orbits in the spherical potential $$\mathrm{\Psi }(r)=\frac{GM}{a}\mathrm{log}\left[\frac{\sqrt{a^2+r^2}+a}{r}\right].$$ (8) The scale length $`a=178.0`$kpc and the mass $`M=2.0\times 10^{12}`$ M are chosen to match those estimated for the halo of the Milky Way (Wilkinson et al. 2003). We investigate whether the trend in Figure 7 can be reproduced by a plausible Galactic orbit as follows. First, we assume that the orbit has a turning point in the RA range $`190200^{}`$ and in the distance range $`5080`$kpc. We then choose values for the line-of-sight distance $`d_0`$, Galactic latitude $`b_0`$ and longitude $`l_0`$ of the turning point and the values of the two components of velocity transverse to the line of sight, $`v_{\mathrm{b},0}`$ and $`v_{\mathrm{l},0}`$. From each set of initial conditions, we integrate the orbit in the RA range $`190200^{}`$ and determine $`V_{\mathrm{gal}}(RA)`$. Assuming Gaussian errors $`\sigma _i`$ on the individual radial velocities (and neglecting any errors in the RA measurements), the probability that the data $`(V_{\mathrm{gal},\mathrm{i}},RA_i)`$ were drawn from the relation $`V_{\mathrm{gal}}(RA)`$ is given by $`P(v_{\mathrm{los},\mathrm{i}},RA_i\|l_0,b_0,r_0,v_{\mathrm{b},0},v_{\mathrm{l},0})=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{\sqrt{2\pi \sigma _i^2}}}\mathrm{exp}\left[{\displaystyle \frac{(v_{\mathrm{los}}(RA_i)v_{\mathrm{los},\mathrm{i}})^2}{2\sigma _i^2}}\right].`$ (9) We use a downhill simplex algorithm (the routine amoeba in Press et al. 1992) to maximise this probability over the five dimensional parameter space. The resulting $`V_{\mathrm{los}}(RA)`$ relation is shown in Figure 7. The orbit we obtain is strongly radial and has an apocentre of $`67.1`$kpc and pericentre of $`7.5`$kpc. The energy and angular momentum of the orbit are $`E=8.1\times 10^4`$(km s$`{}_{}{}^{1})^2`$ and $`L^2=1.2\times 10^7`$(kpc km s$`{}_{}{}^{1})^2`$. The velocity dispersion of all nine stars relative to this orbit is $`15\pm 4`$km s<sup>-1</sup>. Up to this point, we have not made use of the estimated distances to our tracer stars, due to their relatively large uncertainties. We can include them in the orbit determination in a straightforward manner by multiplying the probability in eq (9) by a second Gaussian $`P(d_{\mathrm{los}})`$ given by $`P(d_{\mathrm{los},\mathrm{i}},RA_i\|l_0,b_0,r_0,v_{\mathrm{b},0},v_{\mathrm{l},0})=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \sigma _{d,i}^2}}}\mathrm{exp}\left[{\displaystyle \frac{(d_{\mathrm{los}}(RA_i)d_{\mathrm{los},\mathrm{i}})^2}{2\sigma _{d,i}^2}}\right].`$ (10) Here, $`\sigma _{d,i}`$ is the uncertainty in the line of sight distance to the $`i`$th star. In fact, the inclusion of this term results in an almost identical orbit, the distance uncertainties rendering the distance estimates of little value in the determination of orbital parameters. We initially analysed the data using the technique of Lynden-Bell & Lynden-Bell (1995). These authors note that if a group of stars lie on the same orbit in an assumed spherical potential, they share the same energy $`E`$ and total angular momentum $`L^2`$. Since the energy is given by $$\|E\|=\mathrm{\Psi }(r)\frac{1}{2}v_\mathrm{r}^2\frac{1}{2}\frac{L^2}{r^2},$$ (11) where $`v_\mathrm{r}`$ is the Galactocentric radial velocity, $`L`$ is the magnitude of the total angular momentum vector and $`\mathrm{\Psi }(r)`$ is the potential, the stars in a stream lie on a straight line in a plot of $`\|E_\mathrm{r}\|`$ versus $`r^2`$ where $$\|E_\mathrm{r}\|=\mathrm{\Psi }(r)\frac{1}{2}v_\mathrm{r}^2.$$ (12) If we apply this technique to our data, we obtain orbits with very significantly higher orbital angular momenta than that for the orbit which fits the $`V_{\mathrm{gal}}`$ versus RA relation above. After consideration of the propagation of the observational errors into the $`(\|E_r\|,1/r^2)`$ plane, we noted that the orbital angular momentum was in fact being governed by the orientation of the extended error ellipses caused by the large distance uncertainties. In addition, the derived orbit does not reproduce the trend of $`V_{\mathrm{gal}}`$ with RA seen in Figure 7. We conclude therefore that this technique can only yield useful information about an orbit if the magnitudes of the distance errors are significantly smaller than the radial range covered by the survey – this is not the case for our present sample. Given the proximity of the Sagittarius (Sgr) stream, it is reasonable to suppose the BHB stars might be part of a more distant passage of the stream around the Milky Way (e.g. Helmi & White 2001; Dohm-Palmer et al. 2001). Comparing the positions and heliocentric velocities of our sample of stars to the simulations plotted in Dohm-Palmer et al. (2001; their figure 3) reveal that our sample largely lies either between or beyond two wraps of the Sgr stream. In fact, two of our six BHB stream stars (numbers: 34 and 28) are inconsistent with the model simulations. More recently, Law, Johnston & Majewski (2005) produced models of the Sgr tidal tails using test particle orbits and $`N`$-body simulations in a variety of potentials. Before we are able to compare our data with these models we need to convert our coordinates to the system defined in Majewski et al. (2003). In this coordinate system the zero plane of the latitude coordinate $`B_{}`$ coincides with the best-fit great circle defined by the Sgr debris, as seen from the Sun; the longitudinal coordinate $`\mathrm{\Lambda }_{}`$ is zero in the direction of the Sgr core and increases along the Sgr trailing stream. Our sample resides in the region $`249^{}<\mathrm{\Lambda }_{}<275^{}`$ and $`7^{}<B_{}<20^{}`$. At these coordinates the whole sample of stars in Table 4 resides at larger distances than the models predict. Clearly, the observation of a larger sample of remote BHB stars in the Galaxy halo, along different lines of sight, is essential to confirm the reality of the stream. Additionally, we need to establish whether the small velocity dispersion measured for the eight distant BHB stars discovered in this paper is actually because six of the stars are associated with a coherent structure, or because the velocity dispersion of the whole population of outer halo BHB stars falls steeply with radius. However, we note that the apparent dominance of streaming motion in our BHB sample lends support to the claim of Majewski (2004) that the non-uniform kinematics of outer halo K-giants are consistent with that population having derived almost completely from accretion. We close with a summary of the main points of this paper. We have presented the results of a survey of remote halo A-type stars selected from the SDSS. Spectroscopy of the A-type stars obtained with the VLT produced a sample of 20 stars with data of suitable quality for classification into the classes BHB and A/BS. The final sample (Table 4) comprises eight stars classified BHB, at distances of $`65102`$kpc from the Sun (mean distance $`80`$kpc), with heliocentric radial velocities accurate to 12 km s<sup>-1</sup>, on average, and distance errors $`<10\%`$. This is the most distant sample of Galactic stars with measured radial velocities, of this size. Of the eight remote BHB stars, we find that six show a strong trend in $`V_{\mathrm{gal}}`$ with RA, and are consistent with a single orbit in a spherical halo potential. The measured dispersion of the radial component of the Galactocentric velocity for this sample is $`42\pm 12`$km s<sup>-1</sup>. This value is significantly smaller than values measured for samples of stars at smaller radii, and for satellites at similar radii. This evidence is supported by the existence of three previously identified carbon stars with the same kinematics. A simple model shows all the stars lying on an orbit with energy and angular momentum of $`E=8.1\times 10^4`$(km s$`{}_{}{}^{1})^2`$ and $`L^2=1.2\times 10^7`$(kpc km s$`{}_{}{}^{1})^2`$. The velocity dispersion of the nine stars is 56$`\pm `$13km s<sup>-1</sup>; the dispersion relative to the calculated orbit is 15$`\pm `$4km s<sup>-1</sup>. We conclude that we find a strong indication of the presence of a stream but further observations are required to trace the full extent of this stream on the sky. ## Acknowledgements We thank Kyle Cudworth for providing us with accurate positional information for the radial velocity stars in M5, and J.A.Smith for supplying us with photometric data. We also thank Mike Irwin for several valuable discussions and for pointing us towards the carbon star data set of Totten & Irwin. The Sgr coordinate system conversions made use of code at: http://www.astro.virginia.edu/srm4n/Sgr/code.html. We are grateful to the referee for comments that helped improved the clarity of the manuscript. LC and MIW acknowledge PPARC for financial support. This paper uses observations made on the Very Large Telescope at the European Southern Observatory, Cerro Paranal, Chile \[programme ID: 71.B-0124(A)\]. We made use of the SDSS online database. Funding for the creation and distribution of the SDSS Archive has been provided by the Alfred P. Sloan Foundation, the Participating Institutions, the National Aeronautics and Space Administration, the National Science Foundation, the U.S. Department of Energy, the Japanese Monbukagakusho, and the Max Planck Society.
warning/0506/astro-ph0506286.html	ar5iv	text	# Seismic analysis of the second ionization region of helium in the Sun: I. Sensitivity study and methodology ## 1 Introduction The direct determination of the helium abundance in the solar near-surface layers is difficult and uncertain, although it is very important to the modelling of the internal structure and evolution of the Sun (see Kosovichev et al. 1992 for a comprehensive discussion of the sources of uncertainties). It is usually taken as a fitting parameter of an evolutionary sequence that provides the correct luminosity for the Sun at the present age. The possibility of constraining this parameter to have the observed value for the Sun is important to improve the mass loss estimates and early evolution of the Sun, as well as to test the effects of mixing and settling on stellar evolution. Several attempts have been made to use solar seismic data to calculate the abundance of helium ($`Y`$) in the solar envelope (Dziembowski, Pamyatnykh & Sienkiewicz, 1991; Vorontsov, Baturin & Pamyatnykh, 1991, 1992; Christensen-Dalsgaard & Pérez Hernández, 1992; Pérez Hernández & Christensen-Dalsgaard, 1994; Antia & Basu, 1994; Basu & Antia, 1995; Gough & Vorontsov, 1995; Richard et al., 1998). However the dependence of the determination on other aspects, in particular the equation of state, poses serious difficulties to an accurate direct seismic measurement of the envelope abundance of helium (Kosovichev et al., 1992; Pérez Hernández & Christensen-Dalsgaard, 1994; Basu & Christensen-Dalsgaard, 1997). The sensitivity of the modes to the helium abundance is primarily provided by the change of the local adiabatic sound speed $`c`$ due to ionization. Such sensitivity is given by the behaviour of the first adiabatic exponent, $`\mathrm{\Gamma }_1`$, since $`c^2\mathrm{\Gamma }_1p/\rho `$ where $`p`$ and $`\rho `$ are the pressure and density respectively, and consequently it is strongly dependent on the assumed equation of state and other physics relevant for the region where the ionization takes place. This is the main reason why the seismic determination of the envelope abundance of helium is highly complex. Here we propose a method complementary to those used previously, by adapting the procedure developed by Monteiro et al. (1994, in the following MCDT) and Christensen-Dalsgaard, Monteiro & Thompson (1995, in the following CDMT). In using the solar frequencies in a different way, which provides a direct probe to the region of ionization, we aim at providing a method where the different effects at play in the ionization zone can be isolated, constructing a procedure to access the chemical abundance. Localized variations in the structure of the Sun, such as occur at the base of the convective envelope (see MCDT and Monteiro 1996) and in the region of the second ionization of helium (Monteiro, 1996), create a characteristic signal in the frequencies of oscillation. The properties of such a signal, as measured from the observed frequencies, are related to the location and thermodynamic properties of the Sun at the layer where the sharp or localized variation occurs. The main advantage we see in this method is the possibility to utilise different characteristics of the signal to distinguish different aspects of the physics of the plasma at the region where the signal is generated. In particular we may be able to separate the effects due to convection, the low-temperature opacities and the equation of state from the quantification of the helium abundance that we seek to achieve. Here we mainly concentrate on separating these different contributions in order to establish the dependence of the parameters of the signal in the frequencies on the different aspects of the structure at the ionization region. Using a variational principle we determined how the zone of the second ionization of helium can indeed be considered as a localized perturbation to an otherwise ‘smooth’ structure, generating a characteristic signal in the frequencies of the modes. We note that simplified versions of the expression for the signal discussed here have been applied successfully to cases where there are only very low degree frequencies. The signal has been fitted either to the frequencies of low degree modes (Monteiro & Thompson, 1998; Verner, Chaplin & Elsworth, 2004), or to frequency differences (Miglio et al., 2003; Basu et al., 2004; Vauclair & Théado, 2004; Bazot & Vauclair, 2004; Piau, Ballot & Turck-Chièze, 2005). Here we obtain the expression for the general case of having also modes of higher degree, of which the low degree applications are a particular case. We also demonstrate the method for making use of the information in moderate-degree data available only for the Sun. When using modes with degree above 4 or 5 we can avoid using frequencies affected by the base of the convection zone and may hope to achieve a much higher precision in the results as many more frequencies with lower uncertainties can be used. In this work we present the analysis of the characteristics of the signal under different conditions. Several models with different physics and envelope helium abundances are used to test the method in order to prepare the application to the observed solar data. ## 2 The region of the second ionization of helium In order to model the sensitivity of the modes to this region we must try first to understand how ionization changes the structure. Secondly, we need to estimate how the modes are affected by such a region. The details of the derivations are discussed in the Appendix, but the assumptions and the main expressions are reviewed and analysed here. ### 2.1 Properties of the ionization region Because the helium second ionization zone (Heii ionization zone) is sufficiently deep (well within the oscillatory region of most of the modes) we propose to adapt the method discussed in MCDT to the study of this layer. The contribution from a sharp variation in the structure of the Sun to the frequencies can be estimated by calculating from a variational principle for the modes the effect of a localized feature. In the work by MCDT the feature was the base of the convection zone and the sharp transition was represented by discontinuities in the derivatives of the sound speed. Because of the size of the ionization region when compared with the local wavelength of the modes, that representation is not adequate to reproduce the effect on the frequencies for the ionization region. Here we must, instead, consider how the ionization changes the first adiabatic exponent $`\mathrm{\Gamma }_1(\mathrm{ln}p/\mathrm{ln}\rho )_s`$ (the derivative at constant specific entropy $`s`$) locally, generating what can be described as a ‘bump’ over a region of acoustic thickness of about 300 s (see Fig. 1). This allows us to estimate how the frequencies of oscillation are ‘changed’ due to the presence of this feature in the structure of the Sun. The effect will be mainly taken into account through the changes induced in the adiabatic gradient $`\mathrm{\Gamma }_1`$ by the ionization. Other thermodynamic quantities are also affected, but the changes on the local sound speed is mainly determined by changes in $`\mathrm{\Gamma }_1`$. Therefore, we will calculate the dominant contribution to the changes in the frequencies by establishing what is the effect on the modes due to changes of the adiabatic exponent. Däppen & Gough (1986) and Däppen, Gough & Thompson (1988) have proposed a method based on the same principle, by using the sensitivity of the sound speed to changes on the adiabatic exponent. Using this sensitivity they calibrate a quantity that is associated with ionization in order to try to measure the helium abundance in the solar envelope from seismic data. But most methods have difficulties in removing the dependence of the calibration on the physics of the reference models, making it difficult to obtain a value for the abundance. This is the problem we try to address in this contribution by developing a method able to measure in the frequencies the effect of the ionization and its dependence on the abundance, isolated as much as possible from the other uncertainties. ### 2.2 A variational principle for the effect on the frequencies A variational principle for nonradial adiabatic oscillations, assuming zero pressure at the surface located at radius $`R`$ as a boundary condition, can be formulated. It is possible to further consider only higher-order acoustic modes, for which we may neglect the perturbation in the gravitational potential. The outcome of such a variational principle is an estimate of how the frequencies change due to changes in $`(\mathrm{\Gamma }_1p)`$ and $`\rho `$. This is described and discussed in Appendix A. In order to model the signature of the ionization zone we represent the effect of the second ionization in terms of the changes it induces in the adiabatic exponent $`\mathrm{\Gamma }_1`$. Such a change (see Fig 2) is approximately represented by a ‘bump’ of half width $`\beta `$ in acoustic depth, and relative height $$\delta _\mathrm{d}\left(\frac{\delta \mathrm{\Gamma }_1}{\mathrm{\Gamma }_1}\right)_{\tau _\mathrm{d}},$$ (1) being located at a radial position corresponding to an acoustic depth $`\tau _\mathrm{d}`$. Here, and in the following, acoustic depth $`\tau `$ at a radius $`r`$ is defined as, $$\tau (r)_r^R\frac{\mathrm{d}r}{c},$$ (2) where $`R`$ is the photospheric radius of the Sun. Relatively to the frequencies of a reference model, assumed to be ‘smooth’ and corresponding approximately to a model with no Heii ionization region, we find that the bump changes the frequencies in such a way that there is a periodic component of the form (see Appendix A), $$\delta \omega A(\omega ,l)\mathrm{cos}\mathrm{\Lambda }_\mathrm{d},$$ (3) where the amplitude, as a function of mode frequency $`\omega `$ and mode degree $`l`$, is given by $$A(\omega ,l)a_0\frac{12\mathrm{}/3}{\left(1\mathrm{}\right)^2}\frac{\mathrm{sin}^2\left[\beta \omega \left(1\mathrm{}\right)^{1/2}\right]}{\beta \omega },$$ (4) and the argument is $$\mathrm{\Lambda }_\mathrm{d}2\left[\omega _0^{\tau _\mathrm{d}}\left(1\mathrm{}\right)^{1/2}d\tau +\varphi \right].$$ (5) Here the factor in $`\mathrm{\Delta }`$ represents the geometry of the ray-path accounting for deviation from the vertical when the mode degree is non-zero. It is associated with the Lamb frequency, as given below (Eqs 8, 9). In fact, because the ionization zone is close to the surface and provided we are not using very high-degree data, we can neglect $`\mathrm{}`$ in the expression for the argument $`\mathrm{\Lambda }_\mathrm{d}`$; we can similarly neglect the effect of the mode degree on the surface phase function $`\varphi `$. Consequently, for the ionization zone the expression of the argument becames $$\mathrm{\Lambda }_\mathrm{d}2(\omega \tau _\mathrm{d}+\varphi )2\left(\omega \overline{\tau }_\mathrm{d}+\varphi _0\right).$$ (6) In the asymptotic expression used for the eigenfunction (see Eq. 14), the phase $`\varphi `$ depends on the mode frequency and degree (see MCDT for details). Here we have expanded the phase to first order in frequency, by writing that $`\varphi (\omega )\varphi _0+a_\varphi \omega `$. From this it follows that $`\overline{\tau }_\mathrm{d}\tau _\mathrm{d}+a_\varphi `$, while the frequency independent term of the phase is now $`\varphi _0`$. The amplitude of the signal, through $`a_0`$, corresponds to $$a_0=\frac{3\delta _\mathrm{d}}{2\tau _\mathrm{t}},$$ (7) where $`\tau _\mathrm{t}\tau (0)`$ is the total acoustic size of the Sun. The small factor $`\mathrm{}`$, present in the amplitude, is given by $$\mathrm{}=\mathrm{}_\mathrm{d}\frac{l(l+1)}{\stackrel{~}{l}(\stackrel{~}{l}+1)}\frac{\stackrel{~}{\omega }^2}{\omega ^2},$$ (8) where the value of $`\mathrm{}_\mathrm{d}`$ is given by, $$\mathrm{}_\mathrm{d}=\frac{\stackrel{~}{l}(\stackrel{~}{l}+1)}{\stackrel{~}{\omega }^2}\left(\frac{c}{r}\right)_{\tau =\tau _\mathrm{d}}^2,$$ (9) and $`\stackrel{~}{l}`$ and $`\stackrel{~}{\omega }`$ are two reference values. These values are chosen taking into account the expected behaviour of the signal and the set of modes used, as discussed below. In order to compare the amplitude it is convenient to define a reference value $`A_\mathrm{d}`$, as given by $$A_\mathrm{d}A(\stackrel{~}{\omega },\stackrel{~}{l})=a_0\frac{12\mathrm{}_\mathrm{d}/3}{\left(1\mathrm{}_\mathrm{d}\right)^2}\frac{\mathrm{sin}^2\left[\beta \stackrel{~}{\omega }\left(1\mathrm{}_\mathrm{d}\right)^{1/2}\right]}{\beta \stackrel{~}{\omega }}.$$ (10) The parameters of the signal relevant to characterize the local properties of the ionization zone, as given in Eq. (3), are $`\overline{\tau }_\mathrm{d}`$, $`\beta `$, $`a_0`$ and $`\mathrm{}_\mathrm{d}`$. The values of $`\overline{\tau }_\mathrm{d}`$ and $`\mathrm{}_\mathrm{d}`$ can be used to measure mainly the location of the ionization zone. They both vary strongly with distance to the surface. But the acoustic depth is a cumulative function of the sound speed behaviour over all layers from the surface to a particular position, whereas $`\mathrm{}_\mathrm{d}`$ is a local quantity, not being affected by the behaviour of the sound speed in the layers above the ionzation zone. The values of $`\beta `$ and $`a_0`$ (or $`\delta _\mathrm{d}`$) are expected to be directly related to the local helium abundance, since the size of the bump will be determined by the amount of helium available to be ionized. These parameters are also expected to be strongly affected by the equation of state, and to a lesser extent by the other physics that affect the location of the ionization zone ($`\tau _\mathrm{d}`$). But we may hope to be able to use the other parameters to remove this dependence, while retaining the strong relation between the bump and the helium abundance ($`Y`$). ### 2.3 Measuring the signal in the frequencies Our first goal is to find the five parameters describing the signal from the frequencies of oscillation. In order to do that we must isolate a signature of about 1$`\mu \mathrm{Hz}`$ in amplitude, overimposed in actual frequencies. We do so by isolating in the frequencies the periodic signal described by Eq. (3) using a non-linear least-squares iterative fit to find the best set of parameters. The method used is an adaptation of the one proposed by MCDT; for the present problem we must redefine the characteristic wavelength of the signal to be isolated (quantity $`\lambda _0`$ in MCDT) as it is significantly larger than for the signal from the base of the convective envelope. The parameters describing the signal (Eq. 3), and found by our fitting procedure, are the following; $$\tau _\mathrm{d},\varphi _0,a_0,\mathrm{}_\mathrm{d},\beta .$$ We choose a set of modes which cross the ionization zone, but which do not cross the base of the convection zone. By removing modes that penetrate deep in the Sun (low degree modes), we avoid the contamination coming from the signal generated at the base of the convection zone (see MCDT). But when selecting only modes of higher degree (between 45 and 100), it becomes necessary to include the contribution from the mode degree to the amplitude of the signal. This is the reason why it is necessary to include in the fitting the parameter $`\mathrm{}_\mathrm{d}`$. This parameter is not necessary when studying other stars (Monteiro & Thompson, 1998; Basu et al., 2004; Piau et al., 2005), resulting in a simplified description of the expected observed behaviour. In the case of the Sun there is a great advantage in using all available high-degree modes that cross the ionization zone. The modes considered correspond to the ones available in solar data, having degrees and frequencies such that the lower turning point is between $`0.75R`$ and $`0.95R`$ of the solar radius. The latter ensures the modes cross the ionization zone while the former avoids contamination from the signal originating at the base of the convective envelope (e.g. CDMT, and references therein). These conditions define typically a set of about 450 modes having frequency $`\omega /2\pi `$ in the range $`[1500,3700]\mu \mathrm{Hz}`$, and with mode degree of $`45l100`$. As we are only using modes of high degree in this work, the reference values preferred in the fitting of the signal are; $$\stackrel{~}{l}=100\mathrm{and}\frac{\stackrel{~}{\omega }}{2\pi }=2000\mu \mathrm{Hz}.$$ The first value is an upper limit for modes that cross beyond the ionization zone, while the value of $`\stackrel{~}{\omega }`$ corresponds to the region in frequency where the signal is better defined. These values are only relevant to normalize the parameters fitted for different models. For solar observations only frequencies with a quoted observational error below 0.5 $`\mu \mathrm{Hz}`$ are included. We ensure consistency of the data sets by restricting the selection of mode frequencies from the models to the modes present in the solar data after applying the above selection rules. We stress that the method adopted for removing the smooth component of the frequencies is a key assumption in the process of fitting the signal. Here we use a polynomial fit with a smoothing parameter on the third derivative (see CDMT). In any case, as long as the analyses for different models and for the solar data are consistent, the comparison of the parameters will be independent of the choice on how to describe the smooth component. Such consistency is ensured by using exactly the same set of frequencies and the same numerical parameters for the fitting. ### 2.4 The signal in the solar data To illustrate the signal extraction, the method proposed here was applied to the analysis of solar seismic data from MDI on the SOHO spacecraft (Scherrer et al., 1995). The signal was isolated as described above for the models. The fitted signal of the Sun is shown in Fig. 3a, together with the error bars. In order to illustrate how well the expression for the signal (Eq. 3) fits the data points we also show in Fig. 3b the signal in the frequencies normalized by the amplitude as given in Eq. (4). The quality of the fit done with Eq. (3), confirms the adequacy of the first order analysis developed in Appendix A leading to the expression given by Eq. (4). The values of the parameters found for the data are given in Table 1. From Monte Carlo simulations we have estimated the uncertainty in the determination of the parameters due to observational uncertainties as indicated by the quoted observational errors. The values found, at the 3$`\sigma `$ level, are also listed in Table 1. It is clear that due to the large amplitude of this signature (above $`1\mu \mathrm{Hz}`$ at $`\omega /2\pi =2000\mu \mathrm{Hz}`$) the precision with which the parameters are determined is very high. As long as the method to isolate this characteristic signature is able to remove the “smooth” component, the results can be used with great advantage for testing the zone of the second ionization of helium in the Sun. ## 3 Solar models with different physics In order to establish how sensitive the different characteristics of the signal are to the properties of the ionization zone, and therefore to the aspects of the Sun which affect the ionization zone, we consider different static models of the Sun calculated with the same mass, photospheric radius and luminosity. The profile of the helium abundance in the models is obtained by calibrating with a constant factor a prescribed abundance profile from an evolution model with the age of the Sun (without settling). We note that imposing the same radius and luminosity for all models is the key difference between the analysis presented here and the work by Basu et al. (2004). If the models are not required to have the same luminosity and radius as the Sun, the properties of the ionization zone are not affected in the same way. Consequently the behaviour of the amplitude of the signal in this case is different from what we find when both these conditions are imposed on the models. The aspects of the physics being tested here are the equation of state (EoS), the theory of convection and the opacity. All these aspects affect the ionization zone by changing its location, size and thermodynamic properties. All models were calculated as in Monteiro (1996; see also Monteiro et al. 1996). These are not intended to represent accurately the Sun, but simply to illustrate the usefulness of the method we propose to study a particular region of the solar envelope. As the simplest possible EoS we have used a Saha equation of state with full ionization at high pressure - this corresponds to SEoS in Table 2. As a more complete EoS we have used the CEFF equation of state as described in Christensen-Dalsgaard & Däppen (1992). For the opacities we have considered a simple power law fit (SOp), or the Rosseland mean opacity tables at low temperatures from Kurucz (1991). To include convection we have taken the standard mixing length theory (Böhm-Vitense 1958, parametrized as in Monteiro et al. 1996) or the more recent CGM model (Canuto et al., 1996). As our reference model, in order to illustrate the changes due to the ionization of helium, we have calculated a very simple solar model ($`Z_0`$) with suppressed Heii ionization, by setting the ionization potential to zero. The helium abundance found for each model corresponds to the value that fits the boundary conditions (by scaling a prescribed dependence of the chemical profile, as taken from an evolved solar model). The behaviour of the adiabatic exponent for some of the models (see Table 2), relative to our reference model ($`Z_0`$), is illustrated in Fig. 4. There is a clear difference on the location of the ionization zone ($`\tau _\mathrm{d}`$) when a different EoS is used. The effects of changes in the formulation of convection or in the opacities are much smaller. In order to have models with the same envelope physics, but different helium abundances, we have calculated solar models with the energy generation rate changed by a prescribed factor $`f_ϵ`$ in the emissivity. These are models $`Z_{3l,3h}`$ and $`Z_{5l,5h,5v}`$ which are similar to $`Z_3`$ and $`Z_5`$, respectively, except for the value of $`f_ϵ`$ which is now different from unity. These correspond to models with a different structure of the core but with envelopes with exactly the same set of physics. All differences between these models in the envelope are due to differences in the chemical composition. To illustrate the differences we plot in Fig. 5 the differences in $`\mathrm{\Gamma }_1`$ between models with the same physics but increasing values for the envelope abundance of helium. As the helium abundance increases, there is a corresponding decrease in hydrogen, which results on a slight separation in temperature of the three major ionization regions. Consequently both ionization regions for the helium expand towards higher temperatures. As the bump becomes slightly wider and moves to a higher temperature, the effect on the frequencies is expected to become smaller. For all models we have calculated the frequencies of linear adiabatic oscillations. The set of frequencies for each model, as used to fit the signature of the ionization zone, is described above. The parameters obtained in fitting Eq. (3) to the frequencies of the models listed in Table 2 (excluding $`Z_0`$) are given in Table 3. ## 4 The effect of the physics on the characteristics of the signal The set of solar models considered here, and listed in Table 2, cover three major aspects of the physics which determine the surface structure of the models: the equation of state, the low temperature opacities and the formulation for convection (defining the superadiabatic layer). In order to use the diagnostic potential of this characteristic signature in the frequencies we need to understand how each parameter describing the signal is affected by a specific aspect of the physics defining the structure of the envelope. One would expect that the shape of the bump is strongly determined by the EoS. But any change in the structure that affects the location of the ionization zone will necessarily introduce an effect on the characteristics of the $`\mathrm{\Gamma }_1`$ profile. Consequently we need first to identify the parameters that depend more strongly on the location. These are most likely $`\overline{\tau }_\mathrm{d}`$ and $`\mathrm{}_\mathrm{d}`$. The changes on the upper structure of the envelope are expected to have a direct effect on the turning point of the modes. Consequently we need to look at the parameters that may be affected by the upper reflecting boundary. This is mainly expected to be $`\varphi _0`$. Finally the area of the bump in $`\mathrm{\Gamma }_1`$ in the ionization zone should reflect the local abundance of helium, if the location is well defined. Therefore we will look at $`a_0`$ and $`\beta `$ in order to identify how the helium abundance $`Y`$ defines the characteristics of the signal in the frequencies. ### 4.1 The location of the ionization zone The most easily identifiable characteristic of the signal is its period. This quantity depends strongly on $`\tau _\mathrm{d}`$, but as discussed when writing Eq. (6) the period also contains a contribution from the upper turning point of the modes (where there is a phase shift of the eigenfunction). This means that the period, or precisely $`\overline{\tau }_\mathrm{d}`$, that we measure is not necessarily a good estimate of location $`\tau _\mathrm{d}`$ of the ionization zone. Figure 6a shows the value of $`\overline{\tau }_\mathrm{d}`$, as found from fitting the signal in the frequencies, versus the value of $`\tau _\mathrm{d}`$, as determined from the location of the local minimum of $`\mathrm{\Gamma }_1`$ in the model. There is a difference of up to about 140 s between $`\overline{\tau }_\mathrm{d}`$ and $`\tau _\mathrm{d}`$, and one is not simply a function of the other. The difference between the two comes from $`a_\varphi `$, which measures the leading-order frequency dependence of the phase transition which the eigenfunctions undergo at the upper turning point. This will be strongly affected by the physics that change the structure of the surface, namely convection, EoS, the low temperature opacities, and the structure of the atmosphere. Consequently, we have to use some caution when taking the parameter $`\overline{\tau }_\mathrm{d}`$ from the fit to estimate the location of the ionization region. As an alternative we can consider one of the other parameters which also depends on the position of the ionization zone. This is $`\mathrm{}_\mathrm{d}`$, given in Table 3 for all models and shown in Fig. 6b as a function of the actual acoustic location of the ionization region. The value of $`\mathrm{}_\mathrm{d}`$, defined in Eq. (9), is not sensitive to the layers near the photosphere, as its value is determined exclusively by the sound speed at the ionization zone. However, the determination of this term is associated with a small correction in the amplitude, which makes it more sensitive to the observational errors when fitting the frequencies. Both panels in Fig. 6 show the solar values of $`\overline{\tau }_\mathrm{d}`$ and $`\mathrm{}_\mathrm{d}`$ with 3$`\sigma `$ uncertainties. The values of $`\mathrm{}_\mathrm{d}`$ indicate that all models calculated with the CEFF equation of state give, even if marginally, a location for the ionization zone consistent with the Sun. Finally, the structure at the top of the envelope is also expected to be reflected in the value of $`\varphi _0`$. The value of this parameter for all models is represented in Fig. 7 as a function of the envelope helium abundance. The larger difference is found when changing the EoS (about $`0.06`$). But changes in the opacities also change $`\varphi _0`$ by as much as $`0.01`$, while the theory of convection changes this by about $`0.01`$. It is interesting to confirm that the fitted value of $`\varphi _0`$ is independent of the helium abundance, as one would expect from the analysis leading to the expression of the signal. Consequently $`\varphi _0`$ may allow the separation between the helium abundance and the physics relevant to the outer layers of the Sun because it is insensitive to $`Y`$ whilst being indicative of some near-surface change that may be required in the physics. The solar value for $`\varphi _0`$ is also shown in Fig. 7. Adjustments in the near surface layers seem to be necessary in order to produce models that have a value of $`\varphi _0`$ consistent with the Sun. Changes in the superadiabatic layer or in the surface opacities may be some of the options for reconciling the models with the solar data. ### 4.2 The equation of state From the analysis of the results listed in Table 3, and as discussed in the previous section, the EoS is the most important factor in defining the characteristics of the signal. In Fig. 8 we show the width parameter $`\beta `$ as a function of $`\mathrm{}_\mathrm{d}`$ (a proxy for the location). Models that have the same EoS (CEFF) lie on a common locus in this diagram, as indicated by the dotted line. The position along this line of models all built with the CEFF varies according to changes in the convection or the surface opacities. Models $`Z_1`$ and $`Z_2`$, built with a different EoS, lie in a different region of the diagram. Thus we claim that, with the location of the ionization zone fixed, the width of the bump in $`\mathrm{\Gamma }_1`$ is mainly a function of the EoS, as expected. Consequently, after using $`\varphi _0`$ to test the surface physics, it is possible to combine the constraints provided by $`\mathrm{}_\mathrm{d}`$ and $`\beta `$ to obtain a direct test on the EoS and the location of the ionization zone. Figure 8 also includes the parameters found for the solar data. These are marginally consistent with the expected behaviour found using models calculated with the CEFF equation of state. Other options for the EoS must be considered in an attempt to bring the models closer to the Sun. ### 4.3 The helium abundance in the envelope From the discussion in the previous sections it follows that any determination of the helium abundance requires a careful tuning of the models to the correct structure of the envelope. Such a fine tuning can be performed based on the sensitivity of the eigenfrequencies to the behaviour of the adiabatic exponent in the region where helium undergoes its second ionization. We have found, as discussed above, that: * $`\mathrm{}_\mathrm{d}`$ provides a process to place the ionization zone in the model at the same acoustic depth as for the Sun – this corresponds to adapting mainly the surface layers of the model (atmosphere and/or convection) in order to place the ionization zones at the same acoustic location as measured in the Sun by the solar value of $`\mathrm{}_\mathrm{d}`$; * $`\beta `$ can then be used to adjust the EoS (or more likely to select it from a few candidates) to match the observed behaviour – this corresponds to verifying that the behaviour of $`\beta `$ as a function of the location ($`\mathrm{}_\mathrm{d}`$) in the models includes the observed solar values for these two parameters; * and finally, the parameters $`\overline{\tau }_\mathrm{d}`$ and $`\varphi _0`$ can be combined to adjust the surface physics in the model, in order to recover the observed solar values – this corresponds to adjusting convection (superadiabatic region, mainly), opacities (low temperature range), photosphere, etc, in a complementary way to the first point, until the solar values can be recovered with the models as both parameters are strongly dependent on these aspects of the physics, but quite insensitive to the actual helium abundance. Consequently, we are left with one last parameter, connected with the shape of the bump through $`\delta _\mathrm{d}`$, which is the amplitude of the signal $`a_0`$, or $`A_\mathrm{d}`$. If the model has been adjusted to the observed data using the remaining parameters, then the amplitude will depend mainly on the helium abundance in the model, which can now be compared with the solar abundance. Such a relation provides a measurement of the helium abundance, which complements the boundary condition used in the evolution to fit the model to the present day Sun. Figure 9 illustrates how such a dependence of $`A_\mathrm{d}`$, as defined in Eq. (10), could be constructed after the other aspects of the physics are adjusted. It is worth noting that, as expected from Fig. 5, the amplitude decreases with increasing $`Y`$, since the changes in $`\mathrm{\Gamma }_1`$ due to changes on the hydrogen abundance dominate the variations of the bump. This regime for the inverse dependence of the amplitude of the signal on the abundance of helium is relevant for stars of low effective temperature. That follows from the overlapping of the three ionization zones (Hi, Hei, Heii ). For stars where these are fully separated in temperature it is expected that the amplitude will increase with the abundance of helium. As shown above (see Figs 7 and 8) the models used here are not fully consistent with the physics of the Sun and seem to be only marginally consistent regarding the equation of state that has been used. Consequently the amplitude $`A_\mathrm{d}`$, as found for the solar data, cannot yet be used as an indicator of the helium abundance in the solar envelope. A more adequate calibration of the surface layers in the models must be developed before an estimation for $`Y`$ is inferred from this parameter. The simplified models we are using here to illustrate the applicability of the method have been calculated with scaling a chemical profile determined without including diffusion and settling of helium. This is one of the aspects that needs to be considered in the models in order to move the parameters found for these closer to the solar values. With such a tuning, based on other seismic constraints and on the parameters of the signal discussed here, we have an independent procedure to adjust our models to the Sun in this region near the surface, where the uncertainties in the physics dominate the structure of the models. ## 5 Conclusion In this work we have developed a complementary method to constrain the properties of the helium second ionization region near the surface of the Sun using high degree mode frequencies. The method proposed here can independently test properties of this region, and provides a possible direct measurement of the helium abundance in the envelope. We have shown that some of the parameters characterizing the signature in the frequencies due to this region in the Sun are very sensitive to the EoS used in the calculation of the models, and so can be used to test and constrain the equation of state. Others of the parameters can also provide an important test on the physics affecting the surface regions of the models, namely convection and the low temperature opacities. By combining the diagnostic potential of the five parameters determined from the data with very high precision the helium abundance can be effectively constrained. Here we were mainly concerned to establish the method and demonstrate how it can be used to study the Heii ionization zone in the Sun, and the physics that affect the structure of the Sun in that region. In spite of having used simplified models to represent the Sun we have illustrate the sensitivity of each parameter to the physics, establishing the approach that can be followed when adequate up-to-date evolved solar models are used. Besides the physical ingredients addressed here, aspects like diffusion and settling and improved opacities have to be implemented in order to provide a physically consistent value of the helium abundance. A calibration of the actual solar helium abundance using models with the best up-to-date physics will be the subject of the second paper in this series. ## Acknowledgements We are grateful to S. Basu, J. Christensen-Dalsgaard, M.P. di Mauro and A. Miglio for valuable discussions. This work was supported in part by the Portuguese Fundação para a Ciência e a Tecnologia through projects POCTI/FNU/43658/2001 and POCTI/CTE-AST/57610/2004 from POCTI, with funds from the European programme FEDER. ## Appendix A A variational principle for the Heii ionization zone We consider here a variational principle, following the procedure by Monteiro (1996), for describing how the modes are affected by the presence of the region of the second ionization of the helium. We start by using a variational principle, for small changes of the eigenfrequencies ($`\omega `$) due to small changes of the structure. It can be written (see Christensen-Dalsgaard et al. 1995, and references therein) in the form $$\delta \omega ^2=\frac{\delta I}{I_1}\mathrm{with}I_1\frac{1}{2}\tau _\mathrm{t}E_o^2.$$ (11) Here, $`\tau _\mathrm{t}`$ is the acoustic size of the Sun, and $$\delta I_0^{\tau _t}\left[\left(\delta B_1+\frac{\mathrm{d}\delta B_0}{\mathrm{d}\tau }\right)E_r^2+\delta B_2\frac{\mathrm{d}E_r^2}{\mathrm{d}\tau }+\delta B_3\frac{\mathrm{d}^2E_r^2}{\mathrm{d}\tau ^2}\right]d\tau ,$$ (12) where $`E_r`$ is the normalized radial component of the eigenfunction (with constant amplitude $`E_0`$). The acoustic depth $`\tau `$ is defined in Eq. (2). From asymptotic analysis (see MCDT) we also have that well inside the turning points and for moderate degree modes, $$\frac{\mathrm{d}^2E_r}{\mathrm{d}\tau ^2}\omega ^2\left(1\mathrm{}\right)E_r,$$ (13) or $$E_rE_0\mathrm{cos}\left[\omega _0^\tau \left(1\mathrm{}\right)^{1/2}d\tau +\varphi \right].$$ (14) The changes in the structure relative to the reference (‘smooth’) model are described with the functions $`\delta B_i`$, as given by $$\frac{\delta B_0}{g/c}=\frac{\delta \rho }{\rho },$$ (15) $`{\displaystyle \frac{\delta B_1}{\omega ^2}}`$ $`=`$ $`\{{\displaystyle \frac{1}{1\mathrm{}}}+2\mathrm{}_\rho 2{\displaystyle \frac{13\mathrm{}/2}{(1\mathrm{})^2}}(\mathrm{}_\rho \mathrm{}_\mathrm{c})`$ (18) $`{\displaystyle \frac{1}{(1\mathrm{})^2}}{\displaystyle \frac{(\mathrm{}_\rho \mathrm{}_\mathrm{c})^2}{4\mathrm{}_g}}`$ $`+2\mathrm{}_g{\displaystyle \frac{\mathrm{}(13\mathrm{}/2)}{(1\mathrm{})^2}}\}{\displaystyle \frac{\delta (\mathrm{\Gamma }_1P)}{(\mathrm{\Gamma }_1P)}}+`$ $`+`$ $`\{{\displaystyle \frac{1}{1\mathrm{}}}\mathrm{}_\rho +{\displaystyle \frac{12\mathrm{}}{(1\mathrm{})^2}}(\mathrm{}_\rho \mathrm{}_\mathrm{c})`$ (21) $`+{\displaystyle \frac{\mathrm{}}{(1\mathrm{})^2}}{\displaystyle \frac{(\mathrm{}_\rho \mathrm{}_\mathrm{c})^2}{4\mathrm{}_g}}`$ $`\mathrm{}_g{\displaystyle \frac{\mathrm{}(12\mathrm{})}{(1\mathrm{})^2}}\}{\displaystyle \frac{\delta \rho }{\rho }},`$ $`{\displaystyle \frac{\delta B_2}{g/c}}`$ $`=`$ $`\left[2{\displaystyle \frac{13\mathrm{}/2}{(1\mathrm{})^2}}+{\displaystyle \frac{1\mathrm{}}{2(1\mathrm{})^2}}{\displaystyle \frac{\mathrm{}_\rho \mathrm{}_\mathrm{c}}{2\mathrm{}_g}}\right]{\displaystyle \frac{\delta (\mathrm{\Gamma }_1P)}{(\mathrm{\Gamma }_1P)}}+`$ (22) $`+`$ $`\left[{\displaystyle \frac{12\mathrm{}}{(1\mathrm{})^2}}+{\displaystyle \frac{\mathrm{}}{(1\mathrm{})^2}}{\displaystyle \frac{\mathrm{}_\rho \mathrm{}_\mathrm{c}}{2\mathrm{}_g}}\right]{\displaystyle \frac{\delta \rho }{\rho }},`$ (23) and $$\delta B_3=\frac{1}{2}\frac{1}{1\mathrm{}}\frac{\delta (\mathrm{\Gamma }_1P)}{(\mathrm{\Gamma }_1P)}+\frac{1}{2}\frac{\mathrm{}}{(1\mathrm{})^2}\frac{\delta \rho }{\rho }.$$ (24) where $`r`$, $`\rho `$, $`c`$ and $`g`$ are distance from the centre, density, adiabatic sound speed and gravitational acceleration, respectively. We have also introduced the following quantities $$\mathrm{}=\frac{l(l+1)c^2}{r^2\omega ^2},$$ (25) where $`l`$ is the mode degree, and $`\mathrm{}_\rho `$ $`=`$ $`{\displaystyle \frac{g}{\omega ^2c}}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\tau }}\mathrm{log}\left({\displaystyle \frac{g}{\rho c}}\right),`$ (26) $`\mathrm{}_\mathrm{c}`$ $`=`$ $`{\displaystyle \frac{g}{\omega ^2c}}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\tau }}\mathrm{log}\left({\displaystyle \frac{g}{r^2}}\right),`$ (27) $`\mathrm{}_g`$ $`=`$ $`\left({\displaystyle \frac{g}{\omega c}}\right)^2.`$ (28) These are all first order quantities, compared to unity, because well inside the resonance cavity of the modes the local wavelength is significantly smaller than the scale of variations of the equilibrium quantities. In order to use the expression for $`\delta I`$ from Eq. 12, it is necessary to replace the term in $`(\mathrm{d}\delta B_0/\mathrm{d}\tau )`$ by integrating by parts to obtain for $`\delta I`$; $$\delta I=_{\tau _a}^{\tau _b}[\delta B_1E_r^2+\left(\delta B_2+\delta B_0\right)\frac{\mathrm{d}E_r^2}{\mathrm{d}\tau }+\delta B_3\frac{\mathrm{d}^2E_r^2}{\mathrm{d}\tau ^2}]d\tau .$$ (29) The integration is done only for the region of the ionization zone, starting at $`\tau _a`$ and ending at $`\tau _b`$. Because we are restricting our analysis to localized variations, it is also assumed that the model differences are zero everywhere else. This does not affect our result since we will only take those changes in the frequencies that are not affected by model differences spreading over regions of size of the order of (or larger than) the local wavelength of the modes. We recall, from asymptotic analysis, that $`E_r^2{\displaystyle \frac{E_o^2}{2}}\mathrm{cos}(\mathrm{\Lambda }),`$ $`{\displaystyle \frac{\mathrm{d}E_r^2}{\mathrm{d}\tau }}{\displaystyle \frac{E_0^2}{2}}\mathrm{\hspace{0.33em}2}\omega (1\mathrm{})^{1/2}\mathrm{sin}(\mathrm{\Lambda }),`$ (30) $`{\displaystyle \frac{\mathrm{d}^2E_r^2}{\mathrm{d}\tau ^2}}{\displaystyle \frac{E_0^2}{2}}\mathrm{\hspace{0.33em}4}\omega ^2(1\mathrm{})\mathrm{cos}(\mathrm{\Lambda }).`$ The argument of the trigonometric functions is $$\mathrm{\Lambda }(\tau )2\left[\omega _0^\tau \left(1\mathrm{}\right)^{1/2}d\tau +\varphi \right].$$ (31) After replacing these expressions in the equation for $`\delta I`$, we find $`{\displaystyle \frac{2}{\omega ^2E_o^2}}\delta I`$ $``$ $`{\displaystyle _{\tau _a}^{\tau _b}}\{[{\displaystyle \frac{\delta B_1}{\omega ^2}}4(1\mathrm{})\delta B_3]\mathrm{cos}\mathrm{\Lambda }`$ (33) $`2(1\mathrm{})^{1/2}{\displaystyle \frac{\delta B_2+\delta B_0}{\omega }}\mathrm{sin}\mathrm{\Lambda }\}\mathrm{d}\tau .`$ This expression gives the variational principle for perturbations in the frequencies due to small changes in the structure, as described by $`\delta B_i`$. The next step is to establish what is the effect on the structure of the ionization zone for helium, relative to a model where such a localized effect is not present. In particular, we need to estimate how $`\mathrm{\Gamma }_1`$, $`P`$ and $`\rho `$ are changed from being slowly varying functions of depth to the actual values they have when the second ionization of helium occurs. The difference will produce the $`\delta (\mathrm{\Gamma }_1P)`$ and $`\delta \rho `$ responsible for changing the frequencies, as given in Eqs (15-24). That will allow us to calculate an expression for the characteristic signal we want to isolate in the frequencies. In order to find an expression for the signal we shall first consider that the changes are dominated by $`\mathrm{\Gamma }_1`$. In doing so, we adopt here a different approach from Monteiro (1996), who consider that the dominant contribution could be isolated in the derivative of the sound speed. We do so because the effect of the ionization is better represented as a ‘bump’ in $`\mathrm{\Gamma }_1`$ (see Figs 1 and 2), extending over a localized region of the Sun. Therefore we retain the terms for $`\delta \mathrm{\Gamma }_1`$, and neglect, as a first approximation, the contributions from $`\delta \rho `$ and $`\delta P`$. In doing so we assume that the changes in the sound speed are mainly due to the changes in the adiabatic exponent. Now, relating $`\delta I`$ to the change in the eigenvalue $`\delta \omega `$ (and using Eq. 11) it follows that $`\left[\delta \omega \right]_{\mathrm{\Gamma }_1}`$ $``$ $`{\displaystyle \frac{\left[\delta I\right]_{\mathrm{\Gamma }_1}}{\omega \tau _\mathrm{t}E_o^2}}`$ (34) $``$ $`{\displaystyle \frac{\omega }{2\tau _\mathrm{t}}}{\displaystyle _{\tau _a}^{\tau _b}}\left(f_c\mathrm{cos}\mathrm{\Lambda }+f_s\mathrm{sin}\mathrm{\Lambda }\right){\displaystyle \frac{\delta \mathrm{\Gamma }_1}{\mathrm{\Gamma }_1}}d\tau ,`$ (35) where $`f_s`$ and $`f_c`$ are functions obtained from adding the coefficients of $`\delta \mathrm{\Gamma }_1`$ in the expressions of $`\delta B_0`$, $`\delta B_1`$, $`\delta B_2`$ and $`\delta B_3`$ (see Eq. 33 and Eqs 15-24). At this point we introduce an approximate description of the effect of the second ionization of helium on the adiabatic exponent. As represented in Fig. 2b, we adopt a prescription where the ‘bump’ is approximately described by its half width $`\beta `$ and height $`\delta _\mathrm{d}(\delta \mathrm{\Gamma }_1/\mathrm{\Gamma }_1)_{\tau _\mathrm{d}}`$, with the maximum located at $`\tau _\mathrm{d}`$. This corresponds to considering the following approximating simple expression for $`\delta \mathrm{\Gamma }_1`$: $$\frac{\delta \mathrm{\Gamma }_1}{\mathrm{\Gamma }_1}\delta _\mathrm{d}\{\begin{array}{cc}\left(1+\frac{\tau \tau _d}{\beta }\right)\hfill & \text{}\tau _d(1\alpha )\beta \tau \tau _d\hfill \\ & \\ \left(1\frac{\tau \tau _d}{\beta }\right)\hfill & \text{}\tau _d\tau \tau _d+(1+\alpha )\beta \hfill \\ & \\ 0\hfill & \text{; elsewhere.}\hfill \end{array}$$ (36) The region of the ionization zone starts at $`\tau _a=\tau _\mathrm{d}(1\alpha )\beta `$ and finishes for $`\tau _b=\tau _\mathrm{d}+(1+\alpha )\beta `$, giving that $`\tau _b\tau _a=2\beta `$ is the width. The parameter $`\alpha `$ represents the asymmetry of the bump, and for a first order analysis it does not affect the result. We further consider that the functions $`f_s`$ and $`f_c`$ are slowly varying functions of the structure when compared with the size of the ionization zone $`(2\beta )`$, and so their derivatives can be ignored in the integration. Using this approximation we may integrate Eq. (35) finding that $`\left[\delta \omega \right]_{\mathrm{\Gamma }_1}`$ $``$ $`{\displaystyle \frac{\omega }{2\tau _\mathrm{t}}}\beta \delta _\mathrm{d}\left\{{\displaystyle \frac{\mathrm{sin}\left[\omega \beta (1\mathrm{})^{1/2}\right]}{\omega \beta (1\mathrm{})^{1/2}}}\right\}^2`$ (38) $`\times \left(f_c\mathrm{cos}\mathrm{\Lambda }_\mathrm{d}+f_s\mathrm{sin}\mathrm{\Lambda }_\mathrm{d}\right).`$ All quantities are now evaluated at $`\tau =\tau _\mathrm{d}`$. Taking the dominant contributions (in terms of powers of $`\omega `$ and derivatives of the reference structure – see CDMT for details) of the functions $`f_c`$ and $`f_s`$ (Eq. 33), we can finally write the signal as being $$[\delta \omega ]_{\mathrm{\Gamma }_1}\frac{3\delta _\mathrm{d}}{2\tau _\mathrm{t}}\frac{12\mathrm{}/3}{1\mathrm{}}\frac{\mathrm{sin}^2\left[\omega \beta (1\mathrm{})^{1/2}\right]}{\omega \beta (1\mathrm{})}\mathrm{cos}\mathrm{\Lambda }_\mathrm{d}.$$ (39) This is the expression that describes the ‘additional’ contribution to the frequencies of oscillation $`\omega _{nl}`$ if the region of the second ionization of helium is present. By assuming that we have $$\omega _{nl}[\omega _{nl}]_{\mathrm{smooth}}+[\delta \omega _{nl}]_{\mathrm{\Gamma }_1},$$ (40) it is now possible to try removing the smooth component $`[\omega _{nl}]_{\mathrm{smooth}}`$, by adjusting the frequencies to the expression we have found for the ‘periodic’ component $`[\delta \omega _{nl}]_{\mathrm{\Gamma }_1}`$. In doing so the parameters describing the structure of the Sun at the location $`\tau _\mathrm{d}`$ are determined.
warning/0506/hep-ph0506306.html	ar5iv	text	# Algebraic Approach to Bare Nucleon Matrix Elements of Quark Operators 1footnote 11footnote 1Dedicated to the memory of Professor Gerhard Soff (1949 - 2004). ## I Introduction One ultimate goal of contemporary strong interaction physics is to find a comprehension of the physical properties of hadrons by means of the underlying theory of Quantum Chromodynamics (QCD). Among several methods which provide a link between QCD (quark and gluon) degrees of freedom and the hadronic spectrum are the QCD sum rules which have to be considered as important nonperturbative approach in understanding the physical observables of hadrons. The sum rule method, first developed for the vacuum sumrule , has later been extended to finite density sumrule\_n\_5 ; sumrule\_n\_10 ; sumrule\_n\_15 , finite temperature sumrule\_T\_5 ; sumrule\_T\_10 , and mixed finite density and finite temperature density\_temperature . Within the QCD sum rule approach, and more generally in hadron physics, pions and nucleons have to be considered as important degrees of freedom because the pion is the lightest (Goldstone) meson, while the nucleon is the lightest baryon. In-medium QCD sum rules provide a direct way to relate changes of hadronic properties to changes of the various condensates, i.e. nucleon and pion expectation values of quark and gluon fields. Therefore, expectation values of a local operator $`\widehat{𝒪}`$ taken between these states, $`\pi _{\mathrm{phys}}\|\widehat{𝒪}\|\pi _{\mathrm{phys}}`$ and $`N_{\mathrm{phys}}\|\widehat{𝒪}\|N_{\mathrm{phys}}`$, need to be known. However, the predictive power of the QCD sum rule method in matter meets uncertainties when evaluating condensates, especially higher mass dimension condensates inside the nucleon. Accordingly, the exploration of nucleon matrix elements is presently an active field of hadron physics, cf. matrixlement\_nucleon\_5 ; lattice\_1 . If the operator $`\widehat{𝒪}`$ consists of hadronic fields, then in principle one needs an effective hadronic theory which decribes the interaction between pions and nucleons, respectively, and the hadrons from which the operator $`\widehat{𝒪}`$ is made of for evaluating these matrix elements. However, if one is concerned with pion matrix elements then the use of soft pion theorems lit1 ; lit3 ; lit4 ; lit5 ; hosaka gives in general good estimates for such expressions, which are related to several so called low-energy theorems like Goldberger-Treiman relation Goldberger\_Treiman , Adler-Weisberger sum rule Adler-Weisberger or Cabibbo-Radicati sum rule Cabibbo-Radicati . These soft pion theorems as algebraic tools are based on the hypothesis of partially conserved axial vector current (PCAC) pcac1 ; pcac2 ; pcac3 and postulated current algebra commutation relations lit4 ; lit5 , and allow in general to trace the pion matrix elements of operators made of effective hadronic fields back to vacuum matrix elements. A feature of the soft pion theorems is that they can also be deduced within quark degrees of freedom. Accordingly, pion matrix elements of quark field operators have also been evaluated by means of the soft pion theorem (if we speak about the soft pion theorem then we mean the special theorem considered in the Appendix A which is the relevant one in our context) expressing the pion field and axial vector current, respectively, by interpolating fields made of quark degrees of freedom Hatsuda ; pion . After discovering the powerful method of current algebra for mesons several attempts have been made to investigate the possibilities for extending this algebra to the case of baryons. Especially, the analog hypothesis of a partially conserved baryon current (PCBC) and the related (and postulated) baryon current algebra has been investigated long time ago pcbc1 ; pcbc2 ; pcbc3 ; pcbc4 ; pcbc5 ; pcbc6 . These attempts focussed on the construction of baryon currents by products of nucleon fields. Furthermore, in pcbc7 this procedure has been studied by considering baryon currents made of quark degrees of freedom where several relations between form factors, e.g. baryon-meson vertex form factors, have been obtained. However, it turned out that, while the PCAC directly leads to the mentioned soft pion theorems for evaluating pion matrix elements, the PCBC does not provide a comprehensive algebraic theorem for evaluating nucleon matrix elements. Therefore, up to now for evaluating nucleon matrix elements of an operator $`\widehat{𝒪}`$ consisting of quark fields more involved tools are needed like chiral quark model matrixlement\_nucleon\_5 , lattice evaluations lattice\_1 , or Nambu-Jona-Lasinio model four\_5 . From this point of view it seems very tempting to look for an algebraic approach for evaluating nucleon matrix elements in analogy to the soft pion theorem. Here, by using directly the nucleon field instead the nucleon current, we propose such an algebraic approach for evaluating matrix elements of quark operators taken between a bare nucleon, i.e. the valence quark contribution. To clarify what the terminology ”bare nucleon” means we recall the basic QCD structure of nucleons. From deep inelastic lepton-nucleon scattering (DIS) experiments we know that nucleons are composite color-singlet systems made of partons. In the language of QCD these are three valence quarks with a current quark mass, accompanied by virtual sea quarks and gluons. Accordingly, the physical nucleon state $`\|N_{\mathrm{phys}}`$ is a highly complicated object consisting of many configurations in the Fock space. For instance, in the case of the proton, the Fock expansion begins with the color-singlet state $`\|uud`$ consisting of three valence quarks which is the so called bare proton state, and continues with $`\|uudg`$, $`\|uud\overline{q}q`$ and further sea quark and gluon states that span the degrees of freedom of the proton in QCD. In the low energy region, many properties of the nucleon can rather successfully be described by approximating the virtual sea quarks and gluons by a cloud of mesons, especially pions, surrounding the bare valence quark core. Accordingly, in the pion cloud model, which resembles the Tamm-Dancoff method Tamm\_Dancoff1 ; Tamm\_Dancoff2 ; Tamm\_Dankoff1 ; Tamm\_Dankoff2 ; Tamm\_Dankoff5 ; Tamm\_Dankoff6 the physical nucleon is viewed as a bare nucleon, which accounts for the three valence quarks, accompanied by the pion cloud which accounts for the virtual sea quarks and gluons. Then the Fock representation for the physical nucleon reads Tamm\_Dankoff1 ; Tamm\_Dankoff2 ; Tamm\_Dankoff3 ; tamm1 ; tamm2 ; Z\_Constant $`\|N_{\mathrm{phys}}=Z_N^{1/2}\left(\|N+\varphi _1\|N\pi +\varphi _2\|N\pi \pi +\mathrm{}\right),`$ (1) where the Fock state $`\|N`$ represents a bare nucleon state, $`\|N\pi `$ and $`\|N\pi \pi `$ represent a bare nucleon with one pion and two pions, respectively, and the dots stand for all of the Fock states consisting of one bare nucleon with more than one pion or heavier mesons. The probability amplitudes $`\varphi _n`$ to find the nucleon in the state $`\|Nn\pi `$ can be evaluated by using a Hamiltonian which describes the pion-nucleon interaction Tamm\_Dancoff1 ; Tamm\_Dancoff2 ; Tamm\_Dankoff1 ; Tamm\_Dankoff2 ; Tamm\_Dankoff3 . Then the bare nucleon probability can also be determined and turns out to be $`Z_N0.9`$ Tamm\_Dancoff2 ; tamm1 . Since the deviation of $`Z_N`$ from $`1`$ comes from pion-nucleon interaction one has to put $`Z_N=1`$ if the pion cloud is not taken into account. By using the Fock expansion (1) the expectation value of an observable $`\widehat{𝒪}`$ taken between the physical nucleon states is given by Tamm\_Dancoff2 ; Tamm\_Dankoff1 ; footnote1 , $`N_{\mathrm{phys}}\|\widehat{𝒪}\|N_{\mathrm{phys}}=Z_N(N\|\widehat{𝒪}\|N+`$ $`\varphi _1^2N\pi \|\widehat{𝒪}\|N\pi +\varphi _2^2N\pi \pi \|\widehat{𝒪}\|N\pi \pi +\mathrm{}).`$ (2) The first term on the right side of (2), i.e. the contribution of the bare nucleon without pions, plays an important role for two reasons. First, the bare nucleon is expected to give the main contribution in many cases footnote2 . And second, for the leading chiral correction one needs only the contributions of the lowest-momentum pions in the cloud allowing an application of the soft pion theorem (see Appendix A), which then reduces the pion cloud terms in (2) also to bare nucleon matrix elements krippa1 ; krippa2 . Accordingly, in this paper we focus on bare nucleon matrix elements and propose an algebraic method for evaluating them. This approach seems capable to estimate nucleon matrix elements of quark operators in a straightforward way. We also note that within the algebraic approach new parameters are not necessary since the bare nucleon matrix elements are traced back to vac̃uum matrix elements, like in the soft pion theorem. We apply the method on two-quark, four-quark and, finally, on six-quark operators inside the nucleon which so far have not been evaluated. The paper is organized as follows. In section II we derive an algebraic formula for evaluating matrix elements taken between the state of a bare nucleon. In section III a valence quark field operator with the quantum numbers of a bare nucleon is introduced. A few tests of the nucleon formula on well known bare nucleon matrix elements of two-quark operators are given in section IV.1 (currents) and IV.2 (chiral condensate). In section IV.3 we explore the valence quark contribution of four-quark condensates within the algebraic method developed and assert an interesting agreement with the results of groundstate saturation approximation when taking properly the valence quark contribution. We also compare our findings for the valence quark contribution of four-quark condensates with recently obtained results within a chiral quark model. In section V we evaluate six-quark condensates inside the bare nucleon. A summary of the results and an outlook can be found in section VI. In Appendix A a derivation of the soft pion theorem is given which shows the similarity of it with our algebraic approach. Details of some evaluations are relegated to the Appendix B. ## II Nucleon formula Let $`\widehat{𝒪}(x)`$ be a local operator which may depend on space and time, $`x=(𝐫,t)`$. We are interested in matrix elements taken between two bare nucleon states $`\|N(k,\sigma )`$ with four-momentum $`k`$ and spin $`\sigma `$ (i.e. $`\|N`$ is either a bare proton $`\|p`$ or a bare neutron $`\|n`$ state, which are considered as QCD eigenstates). To derive a formula for such matrix elements between bare nucleons with finite nucleon masses and momenta we first apply the Lehmann-Symanzik-Zimmermann (LSZ) reduction LSZ ; Zuber on one nucleon state, $`N(k_2,\sigma _2)\|\widehat{𝒪}(x)\|N(k_1,\sigma _1)`$ $`=iZ_\mathrm{\Psi }^{1/2}{\displaystyle d^4x_1N(k_2,\sigma _2)\|\mathrm{T}_W\left(\widehat{𝒪}(x)\widehat{\overline{\mathrm{\Psi }}}_N^\alpha (x_1)\right)\|0}`$ $`\times \left(i\gamma _\mu \stackrel{}{}_{x_1}^\mu +M_N\right)_{\alpha \beta }u_N^\beta (k_1,\sigma _1)\mathrm{e}^{ik_1x_1},`$ (3) where the greek letters $`\alpha ,\beta `$ are Dirac indices. The normalization of the nonperturbative QCD vacuum is $`0\|0=1`$, and the normalization for the nucleon state reads $`N(k_2,\sigma _2)\|N(k_1,\sigma _1)=2E_{k_1}(2\pi )^3\delta ^{(3)}(𝐤_1𝐤_2)\delta _{\sigma _1\sigma _2}`$, where $`E_{k_1}=\sqrt{𝐤_1^2+M_N^2}`$. Throughout the paper we take the sum convention: If two Dirac (or later color) indices are equal or not given explicitly, then a sum over them is implied. The four-momenta are on-shell, $`k_1^2=k_2^2=M_N^2`$; for noninteracting nucleons the bare nucleon mass equals the physical nucleon mass, $`M_N=938`$ MeV. The field $`\widehat{\overline{\mathrm{\Psi }}}_N(x_1)`$ is the interacting (adjoint) nucleon field operator. i.e. off-shell. The equal-time anticommutator for the interacting nucleon field operator is the same as for the free fields and reads $`[\widehat{\mathrm{\Psi }}_N^\alpha (𝐫_1,t),\widehat{\mathrm{\Psi }}_N^\beta (𝐫_2,t)]_+=\delta ^{(3)}(𝐫_1𝐫_2)\delta ^{\alpha \beta }.`$ (4) For the wave function renormalization constant we have $`0Z_\mathrm{\Psi }^{1/2}1`$. The free nucleon spinor satisfies $`(\gamma ^\mu k_\mu M_N)u_N(k,\sigma )=0`$, with normalization $`\overline{u}_N(k,\sigma _2)u_N(k,\sigma _1)=2M_N\delta _{\sigma _1\sigma _2}`$. The operator $`\widehat{𝒪}`$ is, for physical reasons, assumed to consist of an even number of fermionic fields, i.e. a bosonic operator, according to which the Wick time-ordering, $`\mathrm{T}_W\widehat{A}(x_1)\widehat{B}(x_2)=\widehat{A}(x_1)\widehat{B}(x_2)\mathrm{\Theta }(t_1t_2)+\widehat{B}(x_2)\widehat{A}(x_1)\mathrm{\Theta }(t_2t_1)`$, has been taken in Eq. (3). We approximate Eq. (3) by introducing a noninteracting nucleon field operator given by $`\widehat{\mathrm{\Psi }}_N^\alpha (x)={\displaystyle }{\displaystyle \frac{d^3𝐤}{(2\pi )^3}}{\displaystyle \frac{1}{2E_k}}{\displaystyle \underset{\sigma =1}{\overset{2}{}}}(\widehat{a}_N(k,\sigma )u_N^\alpha (k,\sigma )\mathrm{e}^{ikx}`$ $`+\widehat{b}_N^{}(k,\sigma )v_N^\alpha (k,\sigma )\mathrm{e}^{ikx}),`$ (5) with the corresponding anticommutator relations in momentum space $`[\widehat{a}_N(k_1,\sigma _1),\widehat{a}_N^{}(k_2,\sigma _2)]_+=[\widehat{b}_N(k_1,\sigma _1),\widehat{b}_N^{}(k_2,\sigma _2)]_+`$ $`=2E_{k_1}(2\pi )^3\delta ^{(3)}(𝐤_1𝐤_2)\delta _{\sigma _1\sigma _2}.`$ (6) Accordingly, $`\|N(k,\sigma )=\widehat{a}_N^{}(k,\sigma )\|0`$. For the noninteracting nucleon field operator $`Z_\mathrm{\Psi }^{1/2}=1`$, and the equation of motion follows from (5), $`\left(i\gamma _\mu \stackrel{}{}_x^\mu M_N\right)\widehat{\mathrm{\Psi }}_N(x)=0`$, and for the adjoint noninteracting nucleon field operator it reads $`\widehat{\overline{\mathrm{\Psi }}}_N(x)\left(i\stackrel{}{}_x^\mu \gamma _\mu +M_N\right)=0`$, respectively. Then one arrives at $`N(k_2,\sigma _2)\|\widehat{𝒪}(x)\|N(k_1,\sigma _1)={\displaystyle d^4x_1\mathrm{e}^{ik_1x_1}\delta (tt_1)}`$ $`\times N(k_2,\sigma _2)\|[\widehat{𝒪}(x),\widehat{\overline{\mathrm{\Psi }}}_N^\alpha (x_1)]_{}\|0\left(\gamma _0\right)_{\alpha \beta }u_N^\beta (k_1,\sigma _1).`$ (7) Applying this procedure on the left nucleon state yields $`N(k_2,\sigma _2)\|\widehat{𝒪}(x)\|N(k_1,\sigma _1)={\displaystyle d^4x_1d^4x_2\mathrm{e}^{ik_1x_1}\mathrm{e}^{ik_2x_2}\delta (tt_1)\delta (tt_2)}`$ $`\times \overline{u}_N^{\beta _2}(k_2,\sigma _2)\left(\gamma _0\right)_{\beta _2\alpha _2}0\|[\widehat{\mathrm{\Psi }}_N^{\alpha _2}(x_2),[\widehat{𝒪}(x),\widehat{\overline{\mathrm{\Psi }}}_N^{\alpha _1}(x_1)]_{}]_+\|0\left(\gamma _0\right)_{\alpha _1\beta _1}u_N^{\beta _1}(k_1,\sigma _1),`$ (8) which is symmetric under the replacement $`0\|[\widehat{\mathrm{\Psi }},[\widehat{𝒪},\widehat{\overline{\mathrm{\Psi }}}]_{}]_+\|00\|[[\widehat{\mathrm{\Psi }},\widehat{𝒪}]_{},\widehat{\overline{\mathrm{\Psi }}}]_+\|0`$. Eq. (8) is the central point of our investigation and we call it nucleon formula. This formula resembles the soft pion theorem given in Appendix A. It is worth to underline that, due to the $`\delta `$-functions in (8), only the equal-time commutator and anticommutator occur. The anticommutator comes into due to the fact that the commutator in (7) between the operator $`\widehat{𝒪}`$ (consisting of an even number of fermionic operators) and the (adjoint) fermionic field operator $`\widehat{\overline{\mathrm{\Psi }}}_N`$ yields an operator consisting of an odd number of fermionic field operators. Therefore, when applying LSZ (cf. Eq. (3)) on the other nucleon state a Dirac time-ordering, $`T_D\widehat{A}(x_1)\widehat{B}(x_2)=\widehat{A}(x_1)\widehat{B}(x_2)\mathrm{\Theta }(t_1t_2)\widehat{B}(x_2)\widehat{A}(x_1)\mathrm{\Theta }(t_2t_1)`$, is needed. The nucleon formula is valid for a noninteracting nucleon with finite mass $`M_N`$ and finite three momentum $`𝐤`$, and in this respect it goes beyond the soft pion theorem, which is valid for pions with vanishing four-momentum only. In the next section we supplement the nucleon formula with a nucleon field operator expressed by quark fields, which allows then the algebraic evaluation of bare nucleon matrix elements of quark operators. We note a remarkable advantage of the algebraic approach. The operator $`\widehat{𝒪}`$ is a composite operator, i.e. a product of field operators taken at the same space-time point. As it stands, such a composite operator needs to be renormalized. Therefore, a renormalization $`\widehat{𝒪}^{\mathrm{ren}}=\widehat{𝒪}\widehat{𝒪}_0`$ (we abbreviate $`\widehat{𝒪}_00\|\widehat{𝒪}\|0`$), which applies for products of noninteracting field operators, has to be implemented Zuber . However, the term $`\widehat{𝒪}_0`$ is a c-number and, according to Eq. (8), does not contribute because of $`[\widehat{𝒪}_0,\widehat{\overline{\mathrm{\Psi }}}_N]_{}=0`$. Another kind of renormalization for products of interacting field operators is also based on subtracting of c-numbers (so called renormalization constants) which vanish when applying the commutator in the nucleon formula. Therefore, one may consider the composite operator $`\widehat{𝒪}(x)`$ in Eq. (8) as a renormalized operator. This feature is also known within PCAC and PCBC algebra, and in particular within the soft pion theorem. ## III Choice of nucleon field operator The nucleon formula (8) can directly be applied on a local operator $`\widehat{𝒪}`$ which consists of bare nucleonic degrees of freedom, e.g. $`\widehat{𝒪}=\widehat{\overline{\mathrm{\Psi }}}_N\widehat{\mathrm{\Psi }}_N,\widehat{\overline{\mathrm{\Psi }}}_N\gamma _\mu \widehat{\mathrm{\Psi }}_N`$, etc. However, we are interested in operators basing on quark degrees of freedom, e.g. $`\widehat{𝒪}=\widehat{\overline{q}}\widehat{q},\widehat{\overline{q}}\gamma _\mu \widehat{q}`$, etc. To make the relation (8) applicable for such cases one needs to decompose the bare nucleon field operator $`\widehat{\mathrm{\Psi }}_N`$ into the three valence quarks, yielding a composite operator $`\widehat{\psi }_N`$ to be specified by now footnote3 . Although there has been considerable success in understanding the properties of nucleons on the basis of their quark substructure as derived within QCD, a rigorous use of QCD for the nucleons is not yet in reach. Therefore, in order to gain a nucleon field operator which shows up the main features (quantum numbers) of the bare nucleon a more phenomenological approach on the basis of the quark-diquark picture of baryons diquarks is used. To be specific we consider the proton. The bare proton state $`\|uud`$ is defined by the $`SU(2)`$ flavor, $`SU(2)`$ spin wavefunction of three valence quarks. In the quark-diquark model of the bare proton two of these valence quarks are regarded as a composite colored particle (diquark) which obeys the Bose statistic and which has a mass of the corresponding meson (e.g. for QCD with $`N_c=2`$ Pauli-Gürsey symmetry Pauli\_Gursey ), i.e. a mass which is significantly larger than the current quark mass. Within quark degrees of freedom the general expression for such a diquark can be written as pcbc7 ; Pisarski $`\widehat{\mathrm{\Phi }}_{q_1q_2}^{\mathrm{a}\mathrm{b}}(x)=\widehat{q}_1^{\mathrm{a}\mathrm{T}}(x)C\mathrm{\Gamma }\widehat{q}_2^\mathrm{b}(x).`$ (9) Here, $`\widehat{q}_1^\mathrm{a}`$ and $`\widehat{q}_2^\mathrm{b}`$ are quark field operators of flavor $`u`$ or $`d`$ with color index a and b, respectively. Throughout the paper all quark field operators are solutions of the full Dirac equation, $`(i\gamma ^\mu \widehat{D}_\mu m_q)\widehat{q}(x)=0`$ with $`\widehat{D}_\mu =_\mu ig_s\widehat{A}_\mu ^a\lambda ^a/2`$, where $`\widehat{A}_\mu ^a`$ are the gluon fields and $`\mathrm{Tr}(\lambda ^a\lambda ^b)=2\delta ^{ab}`$ ($`a,b=1,\mathrm{},8`$ are Gell-Mann indices, which should not be confused with the color indices (in roman style) $`\mathrm{a},\mathrm{b},\mathrm{c}`$ (later also $`\mathrm{i}\mathrm{},\mathrm{n}`$) of quark fields). The equal-time anticommutator for these interacting quark fields is the same as for free quark fields and reads $`[\widehat{q}_\alpha ^\mathrm{a}(𝐫_1,t),\widehat{q}_\beta ^\mathrm{b}(𝐫_2,t)]_+=\delta ^{(3)}(𝐫_1𝐫_2)\delta _{\alpha \beta }\delta ^{\mathrm{a}\mathrm{b}}.`$ (10) The charge conjugation matrix is $`C=i\gamma _0\gamma _2`$, and $`\mathrm{\Gamma }=\{\mathbf{\hspace{0.17em}1}1,\gamma _5,\gamma _\mu ,\gamma _\mu \gamma _5,\sigma _{\mu \nu }\}`$ is an element of the Clifford algebra. $`C`$ changes the parity of $`\mathrm{\Gamma }`$, e.g. $`C\gamma _5`$ has positive parity. The diquark (9), considered as a composite operator made of quark fields, does generally not commute with quark field operators. On the other side, if the diquark is regarded as an effective boson, it commutes with the fermionic quark fields. This feature of the diquark, considered as a bosonic quasiparticle, can be retained on quark level when neglecting the quantum corrections for the quark fields which are participants of the diquark. Accordingly, the diquark is separated into a classical part and a quantum correction $`\widehat{\mathrm{\Phi }}_{q_1q_2}^{\mathrm{a}\mathrm{b}}(x)=q_1^{\mathrm{a}\mathrm{T}}(x)C\mathrm{\Gamma }q_2^\mathrm{b}(x)+\mathrm{\Delta }\widehat{\mathrm{\Phi }}_{q_1q_2}^{\mathrm{a}\mathrm{b}}(x).`$ (11) The classical Dirac spinors $`q_1^\mathrm{a},q_2^\mathrm{b}`$ are solutions of the full Dirac equation $`(i\gamma ^\mu D_\mu m_q)q(x)=0`$. The classical part in (11), $`\mathrm{\Phi }_{q_1q_2}^{\mathrm{a}\mathrm{b}}=q_1^{\mathrm{a}\mathrm{T}}C\mathrm{\Gamma }q_2^\mathrm{b}`$, commutes with quark field operators. To specify the diquark relevant for a proton we note that there are only two structures, $`\mathrm{\Gamma }=\gamma _5`$ and $`\mathrm{\Gamma }=\gamma _5\gamma _0`$, which have positive parity and vanishing total spin, $`J^P=0^+`$ Pisarski . This is in line with interpolating\_1 , where it was found that the proton has indeed a large overlap with the interpolating field $`\widehat{\eta }_p=ϵ^{\mathrm{a}\mathrm{b}\mathrm{c}}\left(\widehat{u}^{\mathrm{a}\mathrm{T}}C\gamma _5\widehat{d}^\mathrm{b}\right)\widehat{u}^\mathrm{c}`$, where $`ϵ^{\mathrm{a}\mathrm{b}\mathrm{c}}`$ is the total antisymmetric tensor. We also remark that in lattice calculations the field $`\widehat{\eta }_p`$ is usually used lattice\_2 ; lattice since this interpolating field has an appropriate nonrelativistic limit. In addition, the field $`\widehat{\eta }_p`$ is also a part of the so called Ioffe interpolating field, which for the proton is given by $`\widehat{\eta }_{\mathrm{Ioffe}}=2ϵ^{\mathrm{a}\mathrm{b}\mathrm{c}}\left((\widehat{u}^{\mathrm{a}\mathrm{T}}C\widehat{d}^\mathrm{b})\gamma _5\widehat{u}^\mathrm{c}(\widehat{u}^{\mathrm{a}\mathrm{T}}C\gamma _5\widehat{d}^\mathrm{b})\widehat{u}^\mathrm{c}\right)`$ ioffe1 . The Ioffe interpolating field is usually used in QCD sum rule evaluations; for a more detailed motivation see also ioffe2 . These properties in mind we take $`\widehat{\eta }_p`$ as a guide for constructing a proton field operator and obtain a semiclassical interpolating proton field by neglecting the quantum correction of the diquark. Further, we assume that any quark of the nucleon can either be a participant of the diquark or can be located outside the diquark. In this line only two different structures for a semiclassical interpolating proton field may occur and, according to this, the general semiclassical field operator for a bare proton is a linear combination of both of them: $`\widehat{\psi }_p^\alpha (x)=ϵ^{\mathrm{a}\mathrm{b}\mathrm{c}}[A_p\left(u^{\mathrm{a}\mathrm{T}}(x)C\gamma _5d^\mathrm{b}(x)\right)\widehat{u}^{\mathrm{c}\alpha }(x)`$ $`+B_p\left(u^{\mathrm{a}\mathrm{T}}(x)C\gamma _5u^\mathrm{b}(x)\right)\widehat{d}^{\mathrm{c}\alpha }(x)].`$ (12) The colorless operator (12) leads to the quantum numbers of a proton (charge, parity, spin, isospin). In the following we evaluate proton matrix elements on quark level by means of the nucleon formula (8) where the field operator $`\widehat{\mathrm{\Psi }}_p(x)`$ is replaced by $`\widehat{\psi }_p(x)`$ given in Eq. (12). Before going further we have to comment on the normalization of the field operator (12), i.e. on the determination of the coefficients $`A_p`$ and $`B_p`$ which, in general, are complex quantities. With the relativistic normalization for the nucleon state and taking into account that there are two u quarks inside the proton we demand numberoperator $`p(k,\sigma _1)\|\widehat{u}_\alpha ^\mathrm{i}(0)\widehat{u}_\alpha ^\mathrm{i}(0)\|p(k,\sigma _2)=4E_k\delta _{\sigma _1\sigma _2}`$. Evaluating this term using the nucleon formula (8) with (12) we find the normalization, at $`x=0`$, to be footnote4 $`\|A_p\|^2ϵ^{\mathrm{abc}}ϵ^{\mathrm{a}^{}\mathrm{b}^{}\mathrm{c}^{}}\delta ^{\mathrm{c}\mathrm{c}^{}}\overline{u}^{\mathrm{a}^{}\mathrm{T}}C\gamma _5\overline{d}^\mathrm{b}^{}u^{\mathrm{a}\mathrm{T}}C\gamma _5d^\mathrm{b}=2.`$ (13) From $`p(k,\sigma _1)\|\widehat{d}_\alpha ^\mathrm{i}(0)\widehat{d}_\alpha ^\mathrm{i}(0)\|p(k,\sigma _2)=2E_k\delta _{\sigma _1\sigma _2}`$ we deduce $`\|B_p\|^2ϵ^{\mathrm{abc}}ϵ^{\mathrm{a}^{}\mathrm{b}^{}\mathrm{c}^{}}\delta ^{\mathrm{c}\mathrm{c}^{}}\overline{u}^{\mathrm{a}^{}\mathrm{T}}C\gamma _5\overline{u}^\mathrm{b}^{}u^{\mathrm{a}\mathrm{T}}C\gamma _5u^\mathrm{b}=1.`$ (14) Formula (8) in combination with the field operator (12) and the normalizations (13) and (14) summarizes our propositions made for obtaining bare proton matrix elements for quark operators. For neutron matrix elements in Eq. (8) we have to insert the semiclassical field operator for the bare neutron, $`\widehat{\psi }_n(x)`$, which is achieved from (12) by interchanging $`ud`$, $`\widehat{u}\widehat{d}`$ and by the replacements $`A_pA_n`$, $`B_pB_n`$. The corresponding normalizations for the bare neutron field operator, i.e. the determination of $`A_n`$ and $`B_n`$, are obtained from Eqs. (13) and (14) by interchanging the up and down quarks, and $`A_pA_n`$, $`B_pB_n`$. ## IV Testing the nucleon formula In the following we will test the outlined formula, Eq. (8) with field operator Eq. (12), and compare with known bare nucleon matrix elements. Throughout the paper we evaluate matrix elements of a composite operator $`\widehat{𝒪}(x)`$ at $`x=0`$ and therefore omit the argument $`x`$ in matrix elements. ### IV.1 Electromagnetic and axial vector current The electromagnetic current for the noninteracting pointlike neutron on hadronic level is zero, due to the vanishing electric charge of the neutron. For the noninteracting pointlike proton it is given by $`\widehat{J}_\mu ^{\mathrm{em}}(x)=e_p\widehat{\overline{\mathrm{\Psi }}}_p(x)\gamma _\mu \widehat{\mathrm{\Psi }}_p(x)`$, where the electric charge of proton equals the elementary electric charge, $`e_p=e`$. Now there are two possibilities to evaluate such a matrix element on hadronic level: either by means of the algebraic approach Eq. (8) and the anticommutator relation (4), or the usual way by means of the field operator Eq. (5) and the anticommutator relations (6). In both cases it is straightforward to show that $`p(k_2,\sigma _2)\|\widehat{J}_\mu ^{\mathrm{em}}\|p(k_1,\sigma _1)=e_p\overline{u}_p(k_2,\sigma _2)\gamma _\mu u_p(k_1,\sigma _1)`$ on effective hadronic level. As a first test of the nucleon formula we verify this relation on quark level where the electromagnetic current is given by $`\widehat{J}_\mu ^{\mathrm{em}}(x)=\frac{2}{3}e\widehat{\overline{u}}(x)\gamma _\mu \widehat{u}(x)\frac{1}{3}e\widehat{\overline{d}}(x)\gamma _\mu \widehat{d}(x)`$. Indeed, by using Eqs. (B) and (LABEL:eq\_application\_10) from Appendix B for the bare proton we get on quark level $`p(k_2,\sigma _2)\|\widehat{J}_\mu ^{\mathrm{em}}\|p(k_1,\sigma _1)=e_p\overline{u}_p(k_2,\sigma _2)\gamma _\mu u_p(k_1,\sigma _1).`$ (15) Similar, for the bare neutron, with Eqs. (78) and (LABEL:eq\_application\_20) from Appendix B, we obtain on quark level $`n(k_2,\sigma _2)\|\widehat{J}_\mu ^{\mathrm{em}}\|n(k_1,\sigma _1)=0`$. Both of these findings are in agreement with the results on effective hadronic level. Now we look at the axial vector current which on hadronic level for a noninteracting pointlike nucleon is given by $`\widehat{A}_\mu ^a(x)=g_A^v\widehat{\overline{\mathrm{\Psi }}}_N(x)\gamma _\mu \gamma _5\frac{\tau ^a}{2}\widehat{\mathrm{\Psi }}_N(x)`$ formfactor , where $`g_A^v`$ is the axial charge of a bare nucleon. For the moment being in this paragraph up to Eq. (16) $`a=1,2,3`$ are isospin indices, and $`\widehat{\mathrm{\Psi }}_N`$, $`\|N`$ and $`u_N`$ are isoduplets. The isospin matrices $`\tau ^a`$ coincide with Pauli’s spin matrices with normalization $`\mathrm{Tr}(\tau ^a\tau ^b)=2\delta ^{ab}`$. Similarly to the case of electromagnetic current, there are two possibilities to evaluate this matrix element on hadronic level: by means of the algebraic approach Eq. (8) and the anticommutator relation (4), or directly by means of the field operator Eq. (5) and the anticommutator relations (6). In both cases one obtains on effective hadronic level the well known result for pointlike nucleons, $`N(k_2,\sigma _2)\|\widehat{A}_\mu ^a\|N(k_1,\sigma _1)=g_A^v\overline{u}_N(k_2,\sigma _2)\gamma _\mu \gamma _5\frac{\tau ^a}{2}u_N(k_1,\sigma _1)`$. We will verify this relation on quark level, where the axial vector current is defined as $`\widehat{A}_\mu ^a(x)=(\widehat{\overline{u}}(x)\widehat{\overline{d}}(x))\gamma _\mu \gamma _5\frac{\tau ^a}{2}(\widehat{u}(x)\widehat{d}(x))^\mathrm{T}`$. To operate with matrix elements between either bare proton states or bare neutron states we use for the nondiagonal cases $`(a=1,2)`$ the assumed isospin symmetry relations, cf. axial\_5 , $`p\|\widehat{\overline{u}}\gamma _\mu \gamma _5\widehat{d}\|n=p\|\widehat{\overline{u}}\gamma _\mu \gamma _5\widehat{u}\widehat{\overline{d}}\gamma _\mu \gamma _5\widehat{d}\|p`$ and $`n\|\widehat{\overline{d}}\gamma _\mu \gamma _5\widehat{u}\|p=n\|\widehat{\overline{d}}\gamma _\mu \gamma _5\widehat{d}\widehat{\overline{u}}\gamma _\mu \gamma _5\widehat{u}\|n`$ footnote5 . Then, taking the solutions of nucleon formula for two-quark operators, Eqs. (B), (LABEL:eq\_application\_10) for proton states, and Eqs. (78), (LABEL:eq\_application\_20) for neutron states (see Appendix B), yields on quark level for the bare (isoduplet) nucleon $`N(k_2,\sigma _2)\|\widehat{A}_\mu ^a\|N(k_1,\sigma _1)`$ $`=\overline{u}_N(k_2,\sigma _2)\gamma _\mu \gamma _5{\displaystyle \frac{\tau ^a}{2}}u_N(k_1,\sigma _1).`$ (16) Comparision of (16) with the result on effective hadronic level yields for the axial charge $`g_A^v=1`$, in fair agreement with the value $`g_A^v0.84`$ deduced from MIT Bag model evaluations and neutron $`\beta `$-decay experiment footnote6 . ### IV.2 Chiral condensate in nucleon The chiral condensate inside the nucleon is related to the pion-nucleon sigma term gasser , $`\sigma _N`$ $`=`$ $`{\displaystyle \frac{m_q}{2M_N}}N_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{u}}\widehat{u}+\widehat{\overline{d}}\widehat{d}\|N_{\mathrm{phys}}(k,\sigma ),`$ where $`2m_q=m_u+m_d`$. A typical value for the pion-nucleon sigma term is $`\sigma _N=45`$ MeV nuclear\_matter\_1 ; nuclear\_matter\_2 . The sigma term can be decomposed, according to Eq. (2), into a valence quark contribution (bare nucleon) and a pion cloud contribution (sea quarks and gluons): $`\sigma _N=\sigma _N^v+\sigma _N^\pi `$. To evaluate $`\sigma _N^v`$ we first consider the u quark chiral condensate inside the bare proton. With Eqs. (B) and (LABEL:eq\_application\_10) one obtains $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\widehat{u}\|p(k_1,\sigma _1)`$ $`=`$ $`2\overline{u}_p(k_2,\sigma _2)u_p(k_1,\sigma _1),`$ $`p(k_2,\sigma _2)\|\widehat{\overline{d}}\widehat{d}\|p(k_1,\sigma _1)`$ $`=`$ $`1\overline{u}_p(k_2,\sigma _2)u_p(k_1,\sigma _1).`$ These relations show the momentum and spin dependence of the chiral condensate inside the bare proton. Of course, for a finite-size nucleon there are additional momentum dependences for which is accounted for by nucleon formfactors. An application of the nucleon formula to the bare neutron reveals the isospin symmetry relations $`n(k_2,\sigma _2)\|\widehat{\overline{u}}\widehat{u}\|n(k_1,\sigma _1)`$ $`=`$ $`p(k_2,\sigma _2)\|\widehat{\overline{d}}\widehat{d}\|p(k_1,\sigma _1),`$ $`n(k_2,\sigma _2)\|\widehat{\overline{d}}\widehat{d}\|n(k_1,\sigma _1)`$ $`=`$ $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\widehat{u}\|p(k_1,\sigma _1).`$ Accordingly, it is only necessary to compare the findings for the proton with results reported in the literature. For the special case $`k_1=k_2`$ and $`\sigma _1=\sigma _2`$ Eq. (LABEL:chiral\_25) simplifies to $`{\displaystyle \frac{1}{2M_N}}p(k,\sigma )\|\widehat{\overline{u}}\widehat{u}\|p(k,\sigma )`$ $`=`$ $`2(2.1),`$ $`{\displaystyle \frac{1}{2M_N}}p(k,\sigma )\|\widehat{\overline{d}}\widehat{d}\|p(k,\sigma )`$ $`=`$ $`1(1.4).`$ (20) The parenthesized values are the findings of Ref. additional\_chiral for the valence quark contribution which well agree with our results. From (20) and isospin symmetry relations (LABEL:chiral\_27) one may now deduce the valence quark contribution to the nucleon sigma term within the algebraic approach, $`\sigma _N^v={\displaystyle \frac{m_q}{2M_N}}N(k,\sigma )\|\widehat{\overline{u}}\widehat{u}+\widehat{\overline{d}}\widehat{d}\|N(k,\sigma )=3m_q.`$ (21) We compare this result with Ref. chiral\_estimate , where the valence quark contribution to the sigma term has been estimated to be $`\sigma _N^v=\sigma _N/(1+G_Sf_\sigma ^2)`$. By using the given values $`G_S=7.91\mathrm{GeV}^2`$ and $`f_\sigma =0.393\mathrm{GeV}`$ one obtains $`\sigma _N^v=20`$ MeV. Accordingly, our result (21) is, for $`m_q7`$ MeV, in good numerical agreement with chiral\_estimate . Finally, by assuming that the contribution of the pion cloud for the physical proton is the same for the chiral u and d quark condensates one can get rid of the term $`\sigma _N^\pi `$ by subtracting the chiral d quark from the chiral u quark condensate. That means the following approximation should be valid $`p_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{u}}\widehat{u}\widehat{\overline{d}}\widehat{d}\|p_{\mathrm{phys}}(k,\sigma )`$ $`p(k,\sigma )\|\widehat{\overline{u}}\widehat{u}\widehat{\overline{d}}\widehat{d}\|p(k,\sigma )=2M_N,`$ (22) where we have used (LABEL:chiral\_25). Indeed, the result (22) is in fair agreement with $`p_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{u}}\widehat{u}\widehat{\overline{d}}\widehat{d}\|p_{\mathrm{phys}}(k,\sigma )=2M_N(M_\mathrm{\Xi }M_\mathrm{\Sigma })/m_s=1.3\mathrm{GeV}`$ obtained in additional\_chiral . ### IV.3 Four-quark condensates Four-quark condensates seem to be quite important in predicting the properties of light vector mesons within the QCD sum rule method zschocke2 . This is related to the fact that in leading-order the chiral condensate is numerically suppressed since it appears in a renormalization invariant contribution $`m_q\widehat{\overline{q}}\widehat{q}`$. Therefore, the gluon condensate and four-quark condensates become numerically more important. However, the numerical values of four-quark condensates are poorly known, and up to now it remains a challenge to estimate their magnitude in a more reliable way. Accordingly, the evaluation of four-quark condensates inside the nucleon is an important issue. Such quantities have been evaluated in nuclear\_matter\_2 within the groundstate saturation approximation, noting the importance of four-quark condensates also for properties of the nucleon within the QCD sum rule approach. An attempt to go beyond the groundstate saturation approximation has been presented in four\_5 , where, by using the Nambu-Jona-Lasinio model and including pions and $`\sigma `$ mesons, correction terms have been obtained. Further evaluations of four-quark condensates beyond the groundstate saturation approximation have been performed in matrixlement\_nucleon\_5 by using a perturbative chiral quark model for describing the nucleons. Later, the results of matrixlement\_nucleon\_5 have been used for evaluating nucleon parameter at finite density within QCD sum rules matrixlement\_nucleon\_7 . In lattice\_1 ; lattice\_2 lattice evaluations for scalar and traceless four-quark operators with non-vanishing twist have been reported. In view of these very few results obtained so far further insight into such condensates is desirable. Before considering this important issue we notice a general decomposition of four-quark condensates. Let $`\widehat{A}`$ and $`\widehat{B}`$ two arbitrary two-quark operators. Then the nucleon expectation value of $`\widehat{𝒪}=\widehat{A}\widehat{B}`$ can be decomposed as footnote8 ; footnote7 $`N_{\mathrm{phys}}\|\widehat{A}\widehat{B}\|N_{\mathrm{phys}}=\widehat{A}_0N_{\mathrm{phys}}\|\widehat{B}\|N_{\mathrm{phys}}`$ $`+N_{\mathrm{phys}}\|\widehat{A}\|N_{\mathrm{phys}}\widehat{B}_0+N_{\mathrm{phys}}\|\widehat{A}\widehat{B}\|N_{\mathrm{phys}}^\mathrm{C},`$ (23) where the first two terms refer to the so called factorization approximation, while the last term is a correction term to the factorization approximation and describes the scattering of a nucleon with $`\widehat{B}`$ into a nucleon and $`\widehat{A}`$, i.e. it is a sum over all connected scattering Feynman diagrams $`N_{\mathrm{phys}}+\widehat{B}N_{\mathrm{phys}}+\widehat{A}`$. The decomposition (23) is matched with the decompositions (24), (25) and(45) given below, as it is seen in footnote7 where we consider an explicit example for the vector channel. The first two terms in (23) scale with $`N_c^2`$ ($`N_c`$ denotes the number of colors, for the moment being taken as a free parameter of QCD), while the correction term scales with $`N_c`$ factorization1 ; factorization2 . Dividing both sides of (23) by $`N_c^2`$ one recognizes that the last term has to be considered as a correction term of the order $`1/N_c`$ factorization1 ; factorization2 . That means that a factorization of four-quark operators in a cold medium is consistent with the large-$`N_c`$ limit factorization1 ; factorization2 , a statement which is also valid in vacuum factorization3 . #### IV.3.1 Flavor-unmixed four-quark condensates We start our investigation with the flavor unmixed four-quark condensates and consider the general expression of two different kinds of flavor unmixed condensates inside the nucleon, namely condensates without and with Gell-Mann matrices $`\lambda ^a`$ ($`a=1,\mathrm{},8`$ are the Gell-Mann indices, which should not be confused with the color indices (in roman style) $`\mathrm{a},\mathrm{b},\mathrm{c}`$ (later also $`\mathrm{i}\mathrm{},\mathrm{n}`$) of quark fields) four\_5 ; nuclear\_matter\_2 ; footnote9 $`N_{\mathrm{phys}}\|\widehat{\overline{q}}\mathrm{\Gamma }_1\widehat{q}\widehat{\overline{q}}\mathrm{\Gamma }_2\widehat{q}\|N_{\mathrm{phys}}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left[\mathrm{Tr}\left(\mathrm{\Gamma }_1\right)\mathrm{Tr}\left(\mathrm{\Gamma }_2\right){\displaystyle \frac{1}{3}}\mathrm{Tr}\left(\mathrm{\Gamma }_1\mathrm{\Gamma }_2\right)\right]\widehat{\overline{q}}\widehat{q}_0N_{\mathrm{phys}}\|\widehat{\overline{q}}\widehat{q}\|N_{\mathrm{phys}}`$ (24) $`+{\displaystyle \frac{1}{16}}\left[\mathrm{Tr}(\mathrm{\Gamma }_1)\mathrm{Tr}(\gamma ^\mu \mathrm{\Gamma }_2)+\mathrm{Tr}(\mathrm{\Gamma }_2)\mathrm{Tr}(\gamma ^\mu \mathrm{\Gamma }_1){\displaystyle \frac{1}{3}}\mathrm{Tr}(\mathrm{\Gamma }_1\gamma ^\mu \mathrm{\Gamma }_2){\displaystyle \frac{1}{3}}\mathrm{Tr}(\mathrm{\Gamma }_2\gamma ^\mu \mathrm{\Gamma }_1)\right]`$ $`\times \widehat{\overline{q}}\widehat{q}_0N_{\mathrm{phys}}\|\widehat{\overline{q}}\gamma _\mu \widehat{q}\|N_{\mathrm{phys}}+N_{\mathrm{phys}}\|\widehat{\overline{q}}\mathrm{\Gamma }_1\widehat{q}\widehat{\overline{q}}\mathrm{\Gamma }_2\widehat{q}\|N_{\mathrm{phys}}^\mathrm{C},`$ $`N_{\mathrm{phys}}\|\widehat{\overline{q}}\mathrm{\Gamma }_1\lambda ^a\widehat{q}\widehat{\overline{q}}\mathrm{\Gamma }_2\lambda ^a\widehat{q}\|N_{\mathrm{phys}}`$ $`=`$ $`{\displaystyle \frac{2}{9}}\mathrm{Tr}\left(\mathrm{\Gamma }_1\mathrm{\Gamma }_2\right)\widehat{\overline{q}}\widehat{q}_0N_{\mathrm{phys}}\|\widehat{\overline{q}}\widehat{q}\|N_{\mathrm{phys}}`$ (25) $`{\displaystyle \frac{1}{9}}\left[\mathrm{Tr}(\mathrm{\Gamma }_1\gamma ^\mu \mathrm{\Gamma }_2)+\mathrm{Tr}(\mathrm{\Gamma }_2\gamma ^\mu \mathrm{\Gamma }_1)\right]\widehat{\overline{q}}\widehat{q}_0N_{\mathrm{phys}}\|\widehat{\overline{q}}\gamma _\mu \widehat{q}\|N_{\mathrm{phys}}+N_{\mathrm{phys}}\|\widehat{\overline{q}}\mathrm{\Gamma }_1\lambda ^a\widehat{q}\widehat{\overline{q}}\mathrm{\Gamma }_2\lambda ^a\widehat{q}\|N_{\mathrm{phys}}^\mathrm{C},`$ where $`\widehat{\overline{q}}\mathrm{}\widehat{q}`$ is either $`\widehat{\overline{u}}\mathrm{}\widehat{u}`$ or $`\widehat{\overline{d}}\mathrm{}\widehat{d}`$ (the dots stand for $`\mathrm{\Gamma }`$ or $`\mathrm{\Gamma }\lambda ^a`$). For the chiral condensate we take $`\widehat{\overline{q}}\widehat{q}_0=(0.250\mathrm{GeV})^3`$. The decompositions of Eqs. (24) and (25) are related to (23) by means of a Fierz rearrangement; an explicit example for the vector channel is given in footnote7 . The last term on the right side of Eqs. (24) and (25) is a correction term four\_5 to the groundstate saturation approximation nuclear\_matter\_2 , describing the scattering process $`N_{\mathrm{phys}}+\widehat{\overline{q}}\mathrm{}\widehat{q}N_{\mathrm{phys}}+\widehat{\overline{q}}\mathrm{}\widehat{q}`$. To get an idea about the magnitude of these correction terms we consider two typical examples. The factorization approximation (24) (i.e. without the correction term) yields for the scalar channel $`N_{\mathrm{phys}}\|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\widehat{q}\|N_{\mathrm{phys}}=0.173\mathrm{GeV}^4`$. In four\_5 the correction term to this groundstate saturation approximation has been found to be $`N_{\mathrm{phys}}\|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\widehat{q}\|N_{\mathrm{phys}}^\mathrm{C}=0.011\mathrm{GeV}^4`$. As another example we consider the vector channel with Gell-Mann matrices. The factorization approximation (25) (i.e. without the correction term) yields $`N_{\mathrm{phys}}\|\widehat{\overline{q}}\gamma _\mu \lambda ^a\widehat{q}\widehat{\overline{q}}\gamma ^\mu \lambda ^a\widehat{q}\|N_{\mathrm{phys}}=0.335\mathrm{GeV}^4`$, while the correction term in four\_5 is $`N_{\mathrm{phys}}\|\widehat{\overline{q}}\gamma _\mu \lambda ^a\widehat{q}\widehat{\overline{q}}\gamma ^\mu \lambda ^a\widehat{q}\|N_{\mathrm{phys}}^\mathrm{C}=0.139\mathrm{GeV}^4`$. Accordingly, the correction to the groundstate saturation approximation in the scalar channel turns out to be less than 10 percent, while in the vector channel with Gell-Mann matrices it is about 30 percent. As we will see, from (24) and (25) the valence quark contribution can be extracted in a unique way. Now we evaluate the valence quark contribution of four-quark condensates, and start to consider the u quark inside the bare proton. Application of the nucleon formula (8) with the composite proton field operator (12) yields $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\|p(k_1,\sigma _1)`$ $`=\overline{u}_p^{\beta _2}(k_2,\sigma _2)(\gamma _0)_{\beta _2\alpha _2}(\gamma _0)_{\alpha _1\beta _1}u_p^{\beta _1}(k_1,\sigma _1)`$ $`\times (\mathrm{\Gamma }_1)^{\alpha \beta }(\mathrm{\Gamma }_2)^{\gamma \delta }{\displaystyle d^3𝐫_1\mathrm{e}^{i𝐤_1𝐫_1}d^3𝐫_2\mathrm{e}^{i𝐤_2𝐫_2}}`$ $`\times 0\|[\widehat{\psi }_p^{\alpha _2}(𝐫_2,0),[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i}\widehat{\overline{u}}_\gamma ^\mathrm{j}\widehat{u}_\delta ^\mathrm{j},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}]_+\|0.`$ (26) By inserting the expression given in Eq. (80) in the Appendix B into (26) one obtains for the bare proton $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\|p(k_1,\sigma _1)`$ $`=`$ $`{\displaystyle \frac{1}{6}}\widehat{\overline{u}}\widehat{u}_0(3\mathrm{Tr}(\mathrm{\Gamma }_1)\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_2u_p(k_1,\sigma _1)+3\mathrm{Tr}(\mathrm{\Gamma }_2)\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1u_p(k_1,\sigma _1)`$ (27) $`\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1\mathrm{\Gamma }_2u_p(k_1,\sigma _1)\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_2\mathrm{\Gamma }_1u_p(k_1,\sigma _1)).`$ In a similar way one obtains for the four-quark condensates involving Gell-Mann matrices $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\lambda ^a\widehat{u}\|p(k_1,\sigma _1)`$ $`={\displaystyle \frac{8}{9}}\widehat{\overline{u}}\widehat{u}_0(\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1\mathrm{\Gamma }_2u_p(k_1,\sigma _1)`$ $`+\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_2\mathrm{\Gamma }_1u_p(k_1,\sigma _1)).`$ (28) For the d flavor we have $`p(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\|p(k_1,\sigma _1)`$ $`={\displaystyle \frac{1}{2}}p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\|p(k_1,\sigma _1),`$ (29) $`p(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\lambda ^a\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\lambda ^a\widehat{d}\|p(k_1,\sigma _1)`$ $`={\displaystyle \frac{1}{2}}p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\lambda ^a\widehat{u}\|p(k_1,\sigma _1).`$ (30) An analog evaluation of these four-quark operators inside the bare neutron gives $`n(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\|n(k_1,\sigma _1)`$ $`=p(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\|p(k_1,\sigma _1),`$ (31) $`n(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\lambda ^a\widehat{u}\|n(k_1,\sigma _1)`$ $`=p(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\lambda ^a\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\lambda ^a\widehat{d}\|p(k_1,\sigma _1),`$ (32) which, like in the case of chiral condensate, reflect the isospin symmetry. By interchanging $`\widehat{u}\widehat{d}`$ on both sides in Eq. (31) and (32) one also gets the d flavor four-quark condensates inside neutron. The results (27) - (32) for the four-quark condensates inside the bare nucleon distinguish between proton and neutron, and they also contain the dependence of these condensates on the momentum of the (pointlike) nucleon. Of course, as in the case of two-quark matrix elements, for a finite-size nucleon there are additional momentum dependences implemented in formfactors. We compare now these findings of Eqs. (27) - (32) with Ref. nuclear\_matter\_2 , i.e. we set $`k_1=k_2`$ and $`\sigma _1=\sigma _2`$, and average over proton and neutron to get the nucleon condensates, $`2N\|\widehat{𝒪}\|N=p\|\widehat{𝒪}\|p+n\|\widehat{𝒪}\|n`$. First we consider the case of scalar four-quark condensate, i.e. $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_2=\mathrm{𝟏}1`$. Then, for the bare nucleon one obtains from Eqs. (27), (29) and (31) $`N(k,\sigma )\|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\widehat{q}\|N(k,\sigma )={\displaystyle \frac{11}{2}}\widehat{\overline{q}}\widehat{q}_0M_N,`$ (33) while, according to the first term in Eq. (24) (the second term vanishes in this special case), for the physical nucleon the result $`N_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\widehat{q}\|N_{\mathrm{phys}}(k,\sigma )={\displaystyle \frac{11}{6}}\widehat{\overline{q}}\widehat{q}_0M_N{\displaystyle \frac{\sigma _N}{m_q}}`$ $`={\displaystyle \frac{11}{6}}\widehat{\overline{q}}\widehat{q}_0M_N\left({\displaystyle \frac{\sigma _N^v}{m_q}}+{\displaystyle \frac{\sigma _N^\pi }{m_q}}\right)`$ (34) has been obtained in nuclear\_matter\_2 . For the last line of Eq. (34) we have used the decomposition $`\sigma _N=\sigma _N^v+\sigma _N^\pi `$, which has already been considered in section IV.2. Comparing both results, Eqs. (33) and (34), one recognizes, by means of relation $`\sigma _N^v/m_q=3`$ (cf. chiral\_estimate and Eq. (21)), that the result (33) is nothing else but just the valence quark contribution of the scalar four-quark condensate inside the nucleon; it is in agreement with the separated valence quark term of Eq. (34). Due to its importance and its instructive property we will also consider the case $`\mathrm{\Gamma }_1=\mathrm{𝟏}1,\mathrm{\Gamma }_2=\gamma _\rho `$. From Eqs. (27), (29) and (31) we find $`N(k,\sigma )\|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\gamma _\rho \widehat{q}\|N(k,\sigma )`$ $`={\displaystyle \frac{5}{4}}\widehat{\overline{q}}\widehat{q}_0\overline{u}_N(k,\sigma )\gamma _\rho u_N(k,\sigma ),`$ (35) which is the contribution of the bare nucleon, i.e. the valence quark contribution. To compare it with the corresponding result of Ref. nuclear\_matter\_2 we first deduce from Eq. (24) that $`N_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\gamma _\rho \widehat{q}\|N_{\mathrm{phys}}(k,\sigma )`$ $`={\displaystyle \frac{5}{6}}\widehat{\overline{q}}\widehat{q}_0N_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{q}}\gamma _\rho \widehat{q}\|N_{\mathrm{phys}}(k,\sigma ).`$ (36) From (36) we have to extract the valence quark contribution by virtue of Eq. (2) (with $`Z_N1`$) $`N_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{q}}\gamma _\rho \widehat{q}\|N_{\mathrm{phys}}(k,\sigma )`$ $`=N(k,\sigma )\|\widehat{\overline{q}}\gamma _\rho \widehat{q}\|N(k,\sigma )+{\displaystyle \underset{n}{}}\varphi _n^2Nn\pi \|\widehat{\overline{q}}\gamma _\rho \widehat{q}\|Nn\pi .`$ The first term on the right side, which is in fact an average over proton and neutron, is the valence quark term we are interested in, while the second term is the pion cloud contribution. With isospin invariance one obtains $`N(k,\sigma )\|\widehat{\overline{q}}\gamma _\rho \widehat{q}\|N(k,\sigma )`$ $`={\displaystyle \frac{1}{2}}\left(p(k,\sigma )\|\widehat{\overline{u}}\gamma _\rho \widehat{u}\|p(k,\sigma )+p(k,\sigma )\|\widehat{\overline{d}}\gamma _\rho \widehat{d}\|p(k,\sigma )\right)`$ $`={\displaystyle \frac{1}{2}}\left(2\overline{u}_p(k,\sigma )\gamma _\rho u_p(k,\sigma )+1\overline{u}_p(k,\sigma )\gamma _\rho u_p(k,\sigma )\right)`$ $`={\displaystyle \frac{3}{2}}\overline{u}_N(k,\sigma )\gamma _\rho u_N(k,\sigma ),`$ (38) where for the last term we have set $`u_p=u_N`$ because of $`M_p=M_N`$. By inserting (IV.3.1) into (36) and using (38) we obtain $`N_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\gamma _\rho \widehat{q}\|N_{\mathrm{phys}}(k,\sigma )`$ $`={\displaystyle \frac{5}{4}}\widehat{\overline{q}}\widehat{q}_0\overline{u}_N(k,\sigma )\gamma _\rho u_N(k,\sigma )`$ $`+{\displaystyle \frac{5}{6}}\widehat{\overline{q}}\widehat{q}_0{\displaystyle \underset{n}{}}\varphi _n^2Nn\pi \|\widehat{\overline{q}}\gamma _\rho \widehat{q}\|Nn\pi .`$ (39) The first term on the right side of Eq. (39) agrees with our result (35), while the second term on the right side of Eq. (39) is a factorization approximation of the expression $`\underset{n}{}\varphi _n^2Nn\pi \|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\gamma _\rho \widehat{q}\|Nn\pi `$. From that it becomes obvious that our result (35) is nothing else but just the valence quark contribution of (39). Other combinations of Clifford matrices, like $`\mathrm{\Gamma }_1=\gamma _5`$ and $`\mathrm{\Gamma }_2=\gamma _\rho `$, with or without Gell-Mann matrices, can be evaluated and compared in the same way footnote10 . This means that within the algebraic approach (8) for evaluating bare nucleon matrix elements we find an agreement for all of the flavor-unmixed four-quark condensates if one takes from the corresponding results of Ref. nuclear\_matter\_2 the valence quark contribution, for instance by means of the decomposition $`\sigma _N=\sigma _N^v+\sigma _N^\pi `$. Therefore, one actually may consider our evaluation as a re-evaluation of the valence quark contribution of the four-quark condensates of Ref. nuclear\_matter\_2 and a confirmation of their results within an independent microscopic approach (quark-diquark picture of the nucleon). But we have to be aware that such an agreement between our algebraic approach and the factorization approximation applies only for the valence quark contribution of nucleon matrix elements. Especially, such an agreement is not expected when taking into account the pion cloud contributions according to Eq. (2). Having found the remarkable agreement with valence quark contribution of the factorization approximation it becomes interesting to compare our results also with other evaluations in the literature. However, it turns out that a comparison with the recent lattice data of Ref. lattice\_1 is quite involved since in lattice\_1 the condensates have been evaluated at a renormalization scale of about $`\mu _{\mathrm{lattice}}^25\mathrm{GeV}^2`$, which is considerably higher than our renormalization point of about $`\mu ^21\mathrm{GeV}^2`$ (our renormalization point is hidden in the chiral condensate, i.e. $`\widehat{\overline{q}}\widehat{q}_0(\mu ^2)`$). To scale the lattice data from $`5\mathrm{GeV}^2`$ down to the hadronic scale of $`1\mathrm{GeV}^2`$ one needs to know the matrix of anomalous dimension for all of the four-quark operators which accounts for the effect of operator mixing among the four-quark condensates. This operator mixing may change considerably the numerical values and even the signs of the four-quark condensates given in lattice\_1 . Another method which seems also capable to compare our results with lattice data would be to scale our renormalization point $`\mu ^2`$ up to the lattice scale $`\mu _{\mathrm{lattice}}^2`$. Such a procedure requires, however, the knowledge of renormalization scale dependence of the chiral condensate. Alltogether, performing these procedures is out of the scope of the present paper and we therefore abandon a comparison of our results with the ones given in Ref. lattice\_1 . In view of the mentioned points and especially in view of an acceptable lucidity of our paper, we restrict ourselves to a comparison with the recently obtained four-quark condensates of Ref. matrixlement\_nucleon\_5 . Due to the specific notation for the four-quark condensates choosen in Ref. matrixlement\_nucleon\_5 we list our results for the valence quark contribution of scalar four-quark condensates for all channels in the same way as in Ref. matrixlement\_nucleon\_5 . Our results for the valence quark contribution can be obtained from Eqs. (27) - (32) by averaging over proton and neutron ($`k_1=k_2,\sigma _1=\sigma _2`$): $`N\|{\displaystyle \frac{2}{3}}\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\widehat{q}{\displaystyle \frac{1}{2}}\widehat{\overline{q}}\lambda ^a\widehat{q}\widehat{\overline{q}}\lambda ^a\widehat{q}\|N`$ $`=0.0733\mathrm{GeV}^4(0.117\mathrm{GeV}^4),`$ (40) $`N\|{\displaystyle \frac{2}{3}}\widehat{\overline{q}}\gamma _5\widehat{q}\widehat{\overline{q}}\gamma _5\widehat{q}{\displaystyle \frac{1}{2}}\widehat{\overline{q}}\gamma _5\lambda ^a\widehat{q}\widehat{\overline{q}}\gamma _5\lambda ^a\widehat{q}\|N`$ $`=0.0147\mathrm{GeV}^4(0.0567\mathrm{GeV}^4),`$ (41) $`N\|{\displaystyle \frac{2}{3}}\widehat{\overline{q}}\gamma _\mu \widehat{q}\widehat{\overline{q}}\gamma ^\mu \widehat{q}{\displaystyle \frac{1}{2}}\widehat{\overline{q}}\gamma _\mu \lambda ^a\widehat{q}\widehat{\overline{q}}\gamma ^\mu \lambda ^a\widehat{q}\|N`$ $`=0.0586\mathrm{GeV}^4(0.0582\mathrm{GeV}^4),`$ (42) $`N\|{\displaystyle \frac{2}{3}}\widehat{\overline{q}}\gamma _\mu \gamma _5\widehat{q}\widehat{\overline{q}}\gamma ^\mu \gamma _5\widehat{q}{\displaystyle \frac{1}{2}}\widehat{\overline{q}}\gamma _\mu \gamma _5\lambda ^a\widehat{q}\widehat{\overline{q}}\gamma ^\mu \gamma _5\lambda ^a\widehat{q}\|N`$ $`=+0.0586\mathrm{GeV}^4(+0.0567\mathrm{GeV}^4),`$ (43) $`N\|{\displaystyle \frac{2}{3}}\widehat{\overline{q}}\sigma _{\mu \nu }\widehat{q}\widehat{\overline{q}}\sigma ^{\mu \nu }\widehat{q}{\displaystyle \frac{1}{2}}\widehat{\overline{q}}\sigma _{\mu \nu }\lambda ^a\widehat{q}\widehat{\overline{q}}\sigma ^{\mu \nu }\lambda ^a\widehat{q}\|N`$ $`=0.176\mathrm{GeV}^4(0.356\mathrm{GeV}^4).`$ (44) The values parenthesized are the results for these condensates as given in Ref. matrixlement\_nucleon\_5 , but for the physical nucleon, i.e. for a nucleon which contains the valence quark, sea quark and gluon contributions footnote11 . Since we have compared our evaluated valence quark contribution with the total contribution for the physical nucleon of Ref. matrixlement\_nucleon\_5 it becomes obvious that one may actually not expect a perfect numerical agreement. The more interesting fact is to notice that the valence quark contribution for the vector and axial vector channel, (42) and (43), respectively, agrees very well with the total contribution for the physical nucleon. For the scalar, axial scalar and tensor channel the signs for the valence quark contribution and total contribution of Ref. matrixlement\_nucleon\_5 agree, while the numerical values differ. This illustrates that the sea quark and gluon contributions are expected to give noticeable contributions. #### IV.3.2 Flavor-mixed four-quark condensates For the flavor-mixed four-quark condensates the general decomposition reads nuclear\_matter\_2 $`N_{\mathrm{phys}}\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\|N_{\mathrm{phys}}`$ $`={\displaystyle \frac{1}{8}}\widehat{\overline{q}}\widehat{q}_0N_{\mathrm{phys}}\|\widehat{\overline{q}}\widehat{q}\|N_{\mathrm{phys}}\mathrm{Tr}\left(\mathrm{\Gamma }_1\right)\mathrm{Tr}\left(\mathrm{\Gamma }_2\right)`$ $`+\widehat{\overline{q}}\widehat{q}_0N_{\mathrm{phys}}\|\widehat{\overline{q}}\gamma _\mu \widehat{q}\|N_{\mathrm{phys}}[\mathrm{Tr}\left(\mathrm{\Gamma }_1\right)\mathrm{Tr}\left(\gamma ^\mu \mathrm{\Gamma }_2\right)`$ $`+\mathrm{Tr}\left(\mathrm{\Gamma }_2\right)\mathrm{Tr}\left(\gamma ^\mu \mathrm{\Gamma }_1\right)]+N_{\mathrm{phys}}\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\|N_{\mathrm{phys}}^\mathrm{C},`$ (45) where the isospin symmetry relations (LABEL:chiral\_27) have been used. The last term in (45) is a correction term to the factorization approximation, describing the scattering process $`N_{\mathrm{phys}}+\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}N_{\mathrm{phys}}+\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}`$. The decomposition (45) is, like (24) and (25), matched with Eq. (23) by means of a Fierz transformation. The flavor-mixed condensates with Gell-Mann matrices vanish in the factorization approximation, $`N_{\mathrm{phys}}\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{d}}\mathrm{\Gamma }_2\lambda ^a\widehat{d}\|N_{\mathrm{phys}}=0`$ nuclear\_matter\_2 (in zschocke1 we have found small corrections to the factorization approximation of flavor-mixed condensates for the vector and axial vector channel with Gell-Mann matrices). Now we consider the valence quark contribution of the flavor mixed four-quark condensates, i.e. the contribution of the bare nucleon. Application of our nucleon formula (8) with (12) yields for the bare proton $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\|p(k_1,\sigma _1)`$ $`={\displaystyle \frac{1}{4}}(2\widehat{\overline{u}}\widehat{u}_0\mathrm{Tr}(\mathrm{\Gamma }_1)\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_2u_p(k_1,\sigma _1)`$ $`+\mathrm{\hspace{0.17em}1}\widehat{\overline{d}}\widehat{d}_0\mathrm{Tr}(\mathrm{\Gamma }_2)\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1u_p(k_1,\sigma _1)),`$ (46) $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{d}}\mathrm{\Gamma }_2\lambda ^a\widehat{d}\|p(k_1,\sigma _1)=0,`$ (47) while for the bare neutron we find $`n(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\|n(k_1,\sigma _1)`$ $`={\displaystyle \frac{1}{4}}(1\widehat{\overline{d}}\widehat{d}_0\mathrm{Tr}(\mathrm{\Gamma }_1)\overline{u}_n(k_2,\sigma _2)\mathrm{\Gamma }_2u_n(k_1,\sigma _1)`$ $`+\mathrm{\hspace{0.17em}2}\widehat{\overline{u}}\widehat{u}_0\mathrm{Tr}(\mathrm{\Gamma }_2)\overline{u}_n(k_2,\sigma _2)\mathrm{\Gamma }_1u_n(k_1,\sigma _1)),`$ (48) $`n(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{d}}\mathrm{\Gamma }_2\lambda ^a\widehat{d}\|n(k_1,\sigma _1)=0.`$ (49) The Eqs. (46) and (48) are generalized expressions of the factorization approximation given in nuclear\_matter\_2 because of the dependence on nucleon momentum and the distinction between proton and neutron. The results of Eqs. (47) and (49) are in agreement with the factorization approximation nuclear\_matter\_2 , but, as mentioned above, beyond the factorization approximation these condensates are nonvanishing zschocke1 . To compare the obtained results of Eqs. (46) and (48) with the factorization approximation, i.e. neglecting the correction term in (45), we consider the special case $`k_1=k_2`$, $`\sigma _1=\sigma _2`$ and $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_2=\mathrm{𝟏}1`$. For the bare nucleon one obtains by averaging over the bare proton and bare neutron, $`N(k,\sigma )\|\widehat{\overline{u}}\widehat{u}\widehat{\overline{d}}\widehat{d}\|N(k,\sigma )=6\widehat{\overline{q}}\widehat{q}_0M_N.`$ (50) This result has to be compared with the corresponding finding of nuclear\_matter\_2 which according to Eq. (45) reads $`N_{\mathrm{phys}}(k,\sigma )\|\widehat{\overline{u}}\widehat{u}\widehat{\overline{d}}\widehat{d}\|N_{\mathrm{phys}}(k,\sigma )=2\widehat{\overline{q}}\widehat{q}_0M_N{\displaystyle \frac{\sigma _N}{m_q}}`$ $`=2\widehat{\overline{q}}\widehat{q}_0M_N\left({\displaystyle \frac{\sigma _N^v}{m_q}}+{\displaystyle \frac{\sigma _N^\pi }{m_q}}\right),`$ (51) where in the last expression we have used the decomposition $`\sigma _N=\sigma _N^v+\sigma _N^\pi `$. By means of the relation $`\sigma _N^v/m_q=3`$ (cf. chiral\_estimate and Eq. (21)) the result (50) is in agreement with the separated valence quark contribution of (51). As in the cases considered in the previous subsection IV.3.1 such an agreement with Ref. nuclear\_matter\_2 can be achieved for all the other combinations of $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ of the Clifford algebra. As in the case of flavor-unmixed four-quark operators we compare our findings for the valence quark contribution of flavor-mixed condensates, (46) and (48), with the total result for the physical nucleon of Ref. matrixlement\_nucleon\_5 to examine the magnitude and sign of our results. According to Eqs. (46) - (49) only the valence quark contribution of the scalar channel does not vanish. Its numerical magnitude $`N\|{\displaystyle \frac{2}{3}}\widehat{\overline{u}}\widehat{u}\widehat{\overline{d}}\widehat{d}{\displaystyle \frac{1}{2}}\widehat{\overline{u}}\lambda ^a\widehat{u}\widehat{\overline{d}}\lambda ^a\widehat{d}\|N`$ $`=0.0586\mathrm{GeV}^4(0.094\mathrm{GeV}^4),`$ (52) turns out to be comparable with the evaluation of matrixlement\_nucleon\_5 for the total contribution of the physical nucleon given in the parentheses in (52). The numerical difference in magnitude is caused by sea quark and gluon contributions. To summarize this section, we have evaluated the valence quark contribution to flavor-unmixed and flavor-mixed four-quark condensates for the u and d flavor inside proton and neutron. The results for the flavor-unmixed operators are given by the Eqs. (27) \- (32), and the results for the flavor-mixed operators are given by the Eqs. (46) - (49). We have seen that our findings for the four-quark condensates within the algebraic approach, which is by far a different method than the used one of Ref. nuclear\_matter\_2 , are in agreement with the large-$`N_c`$ limit factorization1 ; factorization2 and with the results of Ref. nuclear\_matter\_2 when taking from there the valence quark contribution only. It seems admissible, especially in view of the agreement with valence quark contribution of factorization, that the nucleon formula yields reliable results for the valence quark contribution of four-quark condensates inside the nucleon. ## V Six-quark Condensates Six-quark condensates become important mainly for two reasons. First, in the operator product expansion (OPE) of current correlators one usually takes into account all terms up to the order of the four-quark condensates and neglects the contributions of higher order. Such an approximation may work or may not work, depending on the specific physical system under consideration. Accordingly, one has to be aware about the contribution of the next order to decide how good such an approximation is. This is also necessary for the more involved case of finite density, where a Gibbs average over all hadronic states of the correlator under consideration has to be taken. Indeed, a very recent estimate of such higher contributions beyond the four-quark condensate for the nucleon correlator in matter underlines also the importance of an estimate for the six-quark condensates inside the nucleon hungary . And second, it is well known that instantons give rise to corrections to the Wilson coefficients of six-quark condensates six\_quark\_1 ; six\_quark\_2 . These corrections cause a substantial enhancement of the vacuum contribution of six-quark condensates in the OPE of current correlators. To investigate such current correlators at finite density implies the knowledge of the nucleon matrix elements of six-quark condensates. Here, after getting confidence on our proposed approach in the previous sections, we will use the nucleon formula to evaluate the valence quark contribution of six-quark condensates inside the nucleon. Within our algebraic approach by using the nucleon formula (8) with (12) we obtain for the u flavor inside the bare proton $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|p(k_1,\sigma _1)`$ $`=\overline{u}_p^{\beta _2}(k_2,\sigma _2)(\gamma _0)_{\beta _2\alpha _2}(\gamma _0)_{\alpha _1\beta _1}u_p^{\beta _1}(k_1,\sigma _1)`$ $`\times \left(\mathrm{\Gamma }_1\right)^{\alpha \beta }\left(\mathrm{\Gamma }_2\right)^{\gamma \delta }\left(\mathrm{\Gamma }_3\right)^{ϵ\zeta }{\displaystyle d^3𝐫_1\mathrm{e}^{i𝐤_1𝐫_1}d^3𝐫_2\mathrm{e}^{i𝐤_2𝐫_2}}`$ $`0\|[\widehat{\psi }_p^{\alpha _2}(𝐫_2,0),[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\delta ^\mathrm{k}\widehat{\overline{u}}_ϵ^\mathrm{m}\widehat{u}_\zeta ^\mathrm{m},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}]_+\|0.`$ (53) The commutator-anticommutator is given in Eq. (82) in Appendix B for a more general case. According to this result the six-quark condensate inside the bare proton is reduced to a four-quark condensate in vacuum. We note one of these four-quark condensates in vacuum saturation approximation sumrule $`\widehat{u}_\beta ^\mathrm{j}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\delta ^\mathrm{l}\widehat{\overline{u}}_ϵ^\mathrm{m}_0={\displaystyle \frac{1}{(12)^2}}\widehat{\overline{u}}\widehat{u}_0^2`$ $`\times \left(\delta _{\beta \gamma }\delta _{\delta ϵ}\delta ^{\mathrm{j}\mathrm{k}}\delta ^{\mathrm{l}\mathrm{m}}\delta _{\beta ϵ}\delta _{\gamma \delta }\delta ^{\mathrm{j}\mathrm{m}}\delta ^{\mathrm{k}\mathrm{l}}\right).`$ (54) When evaluating all of the four-quark condensates of Eq. (82) in the same way one obtains, by using the normalization (13), $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|p(k_1,\sigma _1)`$ $`=`$ $`{\displaystyle \frac{2}{(12)^2}}\widehat{\overline{u}}\widehat{u}_0^2`$ (55) $`\times [\overline{u}_p(k_2,\sigma _2)(\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{\Gamma }_3+\mathrm{\Gamma }_1\mathrm{\Gamma }_3\mathrm{\Gamma }_2+\mathrm{\Gamma }_2\mathrm{\Gamma }_1\mathrm{\Gamma }_3)u_p(k_1,\sigma _1)+\overline{u}_p(k_2,\sigma _2)(\mathrm{\Gamma }_2\mathrm{\Gamma }_3\mathrm{\Gamma }_1+\mathrm{\Gamma }_3\mathrm{\Gamma }_1\mathrm{\Gamma }_2+\mathrm{\Gamma }_3\mathrm{\Gamma }_2\mathrm{\Gamma }_1)u_p(k_1,\sigma _1)`$ $`3\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1\mathrm{\Gamma }_3u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_2\right)3\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_2\mathrm{\Gamma }_3u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_1\right)3\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1\mathrm{\Gamma }_2u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_3\right)`$ $`3\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_3\mathrm{\Gamma }_2u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_1\right)3\overline{u}_p(k_2,\sigma _2)\left(\mathrm{\Gamma }_2\mathrm{\Gamma }_1\right)u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_3\right)`$ $`3\overline{u}_p(k_2,\sigma _2)\left(\mathrm{\Gamma }_3\mathrm{\Gamma }_1\right)u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_2\right)+9\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_3u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_1\right)\mathrm{Tr}\left(\mathrm{\Gamma }_2\right)`$ $`+9\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_2u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_1\right)\mathrm{Tr}\left(\mathrm{\Gamma }_3\right)+9\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_2\right)\mathrm{Tr}\left(\mathrm{\Gamma }_3\right)`$ $`3\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_3u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_1\mathrm{\Gamma }_2\right)3\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_2u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_1\mathrm{\Gamma }_3\right)3\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1u_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_2\mathrm{\Gamma }_3\right)].`$ For the d flavor six-quark condensate we get $`p(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_3\widehat{d}\|p(k_1,\sigma _1)`$ $`={\displaystyle \frac{1}{2}}p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|p(k_1,\sigma _1).`$ (56) These findings are, to the best of our knowledge, the first attempts to estimate the size of six-quark condensates inside a (bare) nucleon. An analog evaluation for the bare neutron reveals the isospin symmetry relations $`n(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|n(k_1,\sigma _1)`$ $`=p(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_3\widehat{d}\|p(k_1,\sigma _1),`$ (57) $`n(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_3\widehat{d}\|n(k_1,\sigma _1)`$ $`=p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|p(k_1,\sigma _1).`$ (58) Finally, by averaging over the proton matrix elements, Eqs. (55) and (56), and the neutron matrix elements, Eqs. (57) and (58), one gets the six-quark condensates inside a bare nucleon. For instance, the six-quark condensate for the scalar channel is found to be $`N(k,\sigma )\|\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\widehat{q}\widehat{\overline{q}}\widehat{q}\|N(k,\sigma )={\displaystyle \frac{55}{8}}\widehat{\overline{q}}\widehat{q}_0^2M_N.`$ (59) Further, we present results for six-quark condensates which contain Gell-Mann matrices. Note that only nucleon matrix elements of colorless operators do not vanish, i.e. only two Gell-Mann matrices may occur. With the general result (82) in Appendix B and normalization (13) we obtain $`p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|p(k_1,\sigma _1)={\displaystyle \frac{2}{(12)^2}}\widehat{\overline{u}}\widehat{u}_0^2`$ $`\times [{\displaystyle \frac{16}{3}}\overline{u}_p(k_2,\sigma _2)(\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{\Gamma }_3+\mathrm{\Gamma }_1\mathrm{\Gamma }_3\mathrm{\Gamma }_2+\mathrm{\Gamma }_2\mathrm{\Gamma }_1\mathrm{\Gamma }_3)u_p(k_1,\sigma _1)+{\displaystyle \frac{16}{3}}\overline{u}_p(k_2,\sigma _2)(\mathrm{\Gamma }_2\mathrm{\Gamma }_3\mathrm{\Gamma }_1+\mathrm{\Gamma }_3\mathrm{\Gamma }_1\mathrm{\Gamma }_2+\mathrm{\Gamma }_3\mathrm{\Gamma }_2\mathrm{\Gamma }_1)u_p(k_1,\sigma _1)`$ $`16\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_2\mathrm{\Gamma }_1\overline{u}_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_3\right)16\overline{u}_p(k_2,\sigma _2)\mathrm{\Gamma }_1\mathrm{\Gamma }_2\overline{u}_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_3\right)16\overline{u}_p(k_2,\sigma _2)\mathrm{¸}\mathrm{\Gamma }_3\overline{u}_p(k_1,\sigma _1)\mathrm{Tr}\left(\mathrm{\Gamma }_1\mathrm{\Gamma }_2\right)].`$ (60) For the d flavor we get $`p(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\lambda ^a\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\lambda ^a\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_3\widehat{d}\|p(k_1,\sigma _1)`$ $`={\displaystyle \frac{1}{2}}p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|p(k_1,\sigma _1).`$ (61) Finally, evaluating these operators inside the bare neutron we obtain the isospin symmetry relations $`n(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|n(k_1,\sigma _1)`$ $`=p(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\lambda ^a\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\lambda ^a\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_3\widehat{d}\|p(k_1,\sigma _1),`$ (62) $`n(k_2,\sigma _2)\|\widehat{\overline{d}}\mathrm{\Gamma }_1\lambda ^a\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_2\lambda ^a\widehat{d}\widehat{\overline{d}}\mathrm{\Gamma }_3\widehat{d}\|n(k_1,\sigma _1)`$ $`=p(k_2,\sigma _2)\|\widehat{\overline{u}}\mathrm{\Gamma }_1\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_2\lambda ^a\widehat{u}\widehat{\overline{u}}\mathrm{\Gamma }_3\widehat{u}\|p(k_1,\sigma _1).`$ (63) As before, by averaging over the proton matrix elements, Eqs. (60) and (61), and the neutron matrix elements, Eqs. (62) and (63), one obtains the six-quark condensates containing Gell-Mann matrices inside a bare nucleon. The presented findings for six-quark condensates inside the bare nucleon provide basic results for further investigations beyond the order of four-quark condensates within the QCD sum rule approach for the nucleon in matter (for the nucleon sum rule in vacuum see ioffe1 ; ioffe2 , and for the nucleon sum rule in matter see nuclear\_matter\_1 ; nuclear\_matter\_2 ; nuclear\_matter\_3 ). Investigations aiming at predictions beyond the order of four-quark condensates, however, imply in addition to the evaluation of the six-quark condensates also the knowledge of the Wilson coefficients for all of these six-quark condensates. So far, these coefficients in the OPE for the nucleon correlator have been determined up to the order of the four-quark condensates hungary . ## VI Summary An algebraic approach for evaluating bare nucleon matrix elements has been presented. The supposed nucleon formula (8) relates bare nucleon matrix elements to vacuum matrix elements and, therefore, new parameters for evaluating them are not needed. A feature of the algebraic method is that the valence quark contribution of two-quark, four-quark and six-quark condensates can be evaluated on the same footing. One aim of the present paper is to demonstrate how the nucleon formula works and to test it in several cases. In doing so, the nucleon has been considered as a composite pointlike object, described by a valence quark and a valence diquark approximated by two classical Dirac spinors. Neither sea quarks nor gluons, or in a hadronic language no meson cloud, have been taken into account here. Accordingly, the results presented have to be considered as pure valence quark contribution to the matrix elements under consideration. A consideration of the electromagnetic current (15) and the axial vector current (16) for the bare nucleon reveals the expected current structure for a pointlike object. We have evaluated then the valence quark contribution of the chiral condensate (21), finding the relation $`\sigma _N^v=3m_q`$ which turns out to be in numerical agreement with results obtained in an earlier work chiral\_estimate . Furthermore, the valence quark contribution of four-quark condensates has been investigated. Our results are given in Eqs. (27) - (32) for flavor-unmixed operators, and in Eqs. (46) - (49) for flavor-mixed operators. In the special case $`k_1=k_2,\sigma _1=\sigma _2`$ we find an interesting agreement with the groundstate saturation approximation explored in Ref. nuclear\_matter\_2 if one separates the valence quark contribution of four-quark condensates from that results. In this respect our approach yields an independent re-evaluation and confirmation of the results of Ref. nuclear\_matter\_2 , because both methods are different from the conceptional point of view, which is already interesting in itself. Even more, our algebraic approach presented recovers the dependence of condensates on momentum for a pointlike nucleon and distinguishes between proton and neutron matrix elements. In this respect it goes beyond the common factorization approximation. In this context it is worth to underline that the agreement between our algebraic approach and the groundstate saturation approximation has been found for the bare nucleon, and not for the physical nucleon. In Eqs. (40) - (44) and (52) we have compared our results with values of four-quark condensates inside the physical nucleons recently obtained within a chiral quark model matrixlement\_nucleon\_5 . As a further application of nucleon formula we have presented results for six-quark condensates inside the bare nucleon, given in Eqs. (55) - (58) and Eqs. (60) - (63), respectively. These values obtained are, to the best of our knowledge, the first evaluation of six-quark condensates inside (bare) nucleons. Finally, in Eqs. (21), (33) and (59) we have given more explicit examples for the scalar channel of two-quark, four-quark and six-quark operators, respectively, inside the bare nucleon, showing up an interesting alternative change in the algebraic sign from two-quark to four-quark and from four-quark to six-quark condensates. A remark should also be in order about nucleon matrix elements of gluonic operators. Evaluating such operators within the algebraic approach presented requires, in general, the implementation of gluonic degrees of freedom into the composite nucleon field operator (12). Such an implementation might be provided by the quark-gluon interpolating nucleon field discussed in another context in schaefer . The algebraic approach can be extended into several directions. First, the description of the proton core with the field operator (12), and the corresponding one for the bare neutron, can be improved, e.g. by implementing an effective potential between the valence quarks. And second, the pion cloud of nucleon, accounting for virtual sea quarks and gluons inside the physical nucleon, can be implemented within the Tamm-Dancoff method. To get an algebraic approach also for such a case one has to combine the nucleon formula (8) with the soft pion theorem (71). This implies the evaluation of the coefficients $`\varphi _n`$ in (1) within the renormalizable pion-nucleon interaction Hamiltonian, which is therefore a topic of further investigations. Accordingly, for the time being the application of nucleon formula (8) in combination with the field operator (12) has to be considered as a first step in determining more accurately nucleon matrix elements of quark operators. In summary, we arrive at the conclusion that a reliable evaluation of quark operators inside nucleons can be based on a purely algebraic approach. This triggers the hope that predictions of in-medium properties of hadrons become more precise in future. ## VII Acknowledgement The work is supported by BMBF and GSI. We would like to thank Prof. Leonid P. Kaptari, Dr. Hanns-Werner Barz, Dr. Vahtang Gogohia, Dr. Ralf Schützhold, Prof. Vladimir Shabaev, Dr. Gyuri Wolf and Dr. Miklos Zétényi for useful discussions. The authors also thank Prof. Horst Stöcker and Prof. Laszlo P. Csernai for their kind support during the work. One of the authors (S.Z.) thanks for the warm hospitality at the Frankfurt Institute for Advanced Studies (FIAS) in Frankfurt a.M./Germany and at the Research Institute for Particle and Nuclear Physics (KFKI-RMKI) in Budapest/Hungary, and he would also like to express his gratitude to Iris Zschocke and Steffen Köhler for their enduring encouragement during the work. ## Appendix A Soft pion theorem In this Appendix we recall a soft pion theorem relevant for our purposes to show the similarity in derivation and final form of it with the nucleon formula (8). Let us consider the general pion matrix element of an operator $`\widehat{𝒪}(x)`$ which in general may depend on space and time hosaka $`\pi ^b(p_2)\|\widehat{𝒪}(x)\|\pi ^a(p_1)`$ $`=iZ_\phi ^{1/2}{\displaystyle d^4x_1\mathrm{e}^{ip_1x_1}\left(\mathrm{}_{x_1}+m_\pi ^2\right)}`$ $`\times \pi ^b(p_2)\|\mathrm{T}_W\left(\widehat{𝒪}(x)\widehat{\phi }^a(x_1)\right)\|0,`$ (64) where the LSZ reduction formalism has been applied on pion state $`\|\pi ^a(p_1)`$. Here, the letters $`a,b=1,2,3`$ denote isospin indices. The normalization of pion state is $`\pi ^b(p_2)\|\pi ^a(p_1)=2E_{p_1}(2\pi )^3\delta ^{(3)}(𝐩_1𝐩_2)\delta ^{ab}`$, where $`E_{p_1}=\sqrt{𝐩_1^2+m_\pi ^2}`$. The wave function renormalization constant is $`0Z_\phi ^{1/2}1`$. The normalization of nonperturbative QCD vacuum is $`0\|0=1`$. Here, $`x_1=(𝐫_1,t_1)`$ is the space-time four-vector, and $`T_W`$ denotes the Wick time ordering. The states $`\|\pi ^a(p_i)`$ are, of course, on-shell states, i.e. solutions of the Klein-Gordon equation for noninteracting pions, while the field operator $`\widehat{\phi }^a`$, in general, is the interacting field, i.e. it is off-shell. The four momenta in (64) are on-shell, i.e. $`p_1^2=p_2^2=m_\pi ^2`$. The soft pion theorem is valid for a noninteracting pion field (i.e. $`Z_\phi ^{1/2}=1`$), wich is a solution of the Klein-Gordon equation $`\left(\mathrm{}_x+m_\pi ^2\right)\widehat{\phi }^a(x)`$. In order to be complete in the representation we will also specify the noninteracting pion field operator $`\widehat{\phi }^a(x)={\displaystyle \frac{d^3𝐩}{(2\pi )^3}\frac{1}{2E_p}\left(\widehat{a}^a(p)\mathrm{e}^{ipx}+\widehat{b}^a(p)\mathrm{e}^{ipx}\right)},`$ (65) where the creation and annihilation operators obey the following commutator relations $`[\widehat{a}^a(p_1),\widehat{a}^b(p_2)]_{}=[\widehat{b}^a(p_1),\widehat{b}^b(p_2)]_{}`$ $`=2E_{p_1}(2\pi )^3\delta ^{(3)}(𝐩_1𝐩_2)\delta ^{ab}.`$ (66) Accordingly, $`\|\pi ^a(p)=\widehat{a}^a(p)\|0`$. From (65) and (66) we deduce the equal-time commutator for the noninteracting pion fields, $`[\widehat{\phi }^a(𝐫_1,t),_0\widehat{\phi }^b(𝐫_2,t)]_{}=i\delta ^{(3)}(𝐫_1𝐫_2)\delta ^{ab}.`$ (67) In addition, the soft pion theorem is only valid in case of vanishing four-momentum $`p_1^\mu 0`$ which implies $`m_\pi =0`$. Then we get $`\underset{\genfrac{}{}{0pt}{}{}{p_10}}{lim}\pi ^b(p_2)\|\widehat{𝒪}(x)\|\pi ^a(p_1)`$ $`=i{\displaystyle d^4x_1\mathrm{}_{x_1}\pi ^b(p_2)\|\mathrm{T}_W\left(\widehat{𝒪}(x)\widehat{\phi }^a(x_1)\right)\|0}`$ $`=i\pi ^b(p_2)\|[\widehat{𝒪}(x),{\displaystyle \frac{}{t_1}}\widehat{\phi }^a(x_1)]_{}\|0\delta (tt_1).`$ In the last line we have used the equation of motion for the massless pion field, i.e. $`\mathrm{}_{x_1}\widehat{\phi }^a(x_1)=0`$. Now, the PCAC hypothesis, which asserts a relation between the axial current and the pion field ($`f_\pi 92.4`$ MeV is the pion decay constant), $`\widehat{A}_\mu ^a(x)=f_\pi _\mu \widehat{\phi }^a(x),`$ (69) is inserted into Eq. (A) (by means of the field equation for the noninteracting pion it becomes obvious from (69) that in the limit of vanishing pion mass the PCAC goes over to a conserved axial vector current (see also formfactor )). The axial vector current $`\widehat{A}_\mu ^a(x)`$ obeys the well known current algebra commutation relations lit4 which directly leads to the soft pion theorem relevant for our purposes hosaka $`\underset{\genfrac{}{}{0pt}{}{}{p_10}}{lim}\pi ^b(p_2)\|\widehat{𝒪}(x)\|\pi ^a(p_1)={\displaystyle \frac{i}{f_\pi }}{\displaystyle d^4x_1}`$ $`\times \pi ^b(p_2)\|[\widehat{𝒪}(x),\widehat{A}_0^a(x_1)]_{}\|0\delta (tt_1).`$ (70) One may apply the same steps as before on the other pion state as well, ending up with the soft pion theorem $`\underset{\genfrac{}{}{0pt}{}{}{p_20}}{lim}\underset{\genfrac{}{}{0pt}{}{}{p_10}}{lim}\pi ^b(p_2)\|\widehat{𝒪}(x)\|\pi ^a(p_1)`$ $`={\displaystyle \frac{1}{f_\pi ^2}}{\displaystyle d^4x_1d^4x_2\delta (tt_1)\delta (tt_2)}`$ $`\times 0\|[\widehat{A}_0^b(x_2),[\widehat{𝒪}(x),\widehat{A}_0^a(x_1)]_{}]_{}\|0`$ $`={\displaystyle \frac{1}{f_\pi ^2}}0\|[\widehat{Q}_A^b,[\widehat{𝒪}(x),\widehat{Q}_A^a]_{}]_{}\|0.`$ (71) where $`\widehat{Q}_A^a`$ is the (time independent) axial charge, $`\widehat{Q}_A^a=d^3𝐫\widehat{A}_0^a(𝐫,t)`$. It seems expedient to emphasize that the QCD quark degrees of freedom were not necessary for deriving the soft pion theorem. If one expresses the axial vector current by quark fields, $`\widehat{A}_\mu ^a(x)=\widehat{\overline{\mathrm{\Psi }}}(\tau ^a/2)\gamma _\mu \gamma _5\widehat{\mathrm{\Psi }}`$ with $`\widehat{\mathrm{\Psi }}=(\widehat{u}\widehat{d})^\mathrm{T}`$, and the pion fields as well by means of their interpolating fields (i.e. composite quark fields which have the quantum numbers of pions), then the relation (69), the current algebra and, therefore, the soft pion theorem (71) can also be established within QCD degrees of freedom. This theorem can then also be used for evaluating pion matrix elements of quark operators (see, for instance, Hatsuda ; pion ). Summarizing, the soft pion theorem is valid for a noninteracting pion field with vanishing pion four-momentum. In many applications such a restriction is not problematic since the pion mass is small compared to a typical hadronic scale of about 1 GeV. ## Appendix B Evaluating Eq. (53) To evaluate Eq. (53) we start with the case of two quark field operators. The nucleon formula (8) with (12) yields $`p(k_2,\sigma _2)\|\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i}\|p(k_1,\sigma _1)`$ $`=\overline{u}_p^{\beta _2}(k_2,\sigma _2)(\gamma _0)_{\beta _2\alpha _2}(\gamma _0)_{\alpha _1\beta _1}u_p^{\beta _1}(k_1,\sigma _1)`$ $`\times {\displaystyle }d^3𝐫_1{\displaystyle }d^3𝐫_2\mathrm{e}^{i𝐤_1𝐫_1}\mathrm{e}^{i𝐤_2𝐫_2}`$ $`\times 0\|[\widehat{\psi }_p^{\alpha _2}(𝐫_2,0),[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}]_+\|0.`$ (72) Inserting the proton field operator (12) and using $`[\widehat{A}\widehat{B},\widehat{C}]_{}=\widehat{A}[\widehat{B},\widehat{C}]_+[\widehat{A},\widehat{C}]_+\widehat{B}`$ for the commutator we obtain $`[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}`$ $`=`$ $`A_p^{}ϵ^{\mathrm{a}\mathrm{b}\mathrm{c}}\left(\overline{u}^{\mathrm{a}\mathrm{T}}(𝐫_1)C\gamma _5\overline{d}^\mathrm{b}(𝐫_1)\right)`$ (73) $`\times \delta ^{\mathrm{i}\mathrm{c}}(\gamma _0)_{\beta \alpha _1}\delta ^{(3)}(𝐫_1)\widehat{\overline{u}}_\alpha ^\mathrm{i}.`$ In the same way we get $`[\widehat{\psi }_p^{\alpha _2}(𝐫_2,0),[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}]_+`$ $`=\|A_p\|^2ϵ^{\mathrm{a}\mathrm{b}\mathrm{c}}ϵ^{\mathrm{a}^{}\mathrm{b}^{}\mathrm{c}^{}}(\gamma _0)_{\beta \alpha _1}(\gamma _0)_{\alpha \alpha _2}`$ $`\times \delta ^{\mathrm{c}\mathrm{c}^{}}\delta ^{(3)}(𝐫_1)\delta ^{(3)}(𝐫_2)`$ $`\times \left(u^{\mathrm{a}^{}\mathrm{T}}(𝐫_2)C\gamma _5d^\mathrm{b}^{}(𝐫_2)\right)\left(\overline{u}^{\mathrm{a}\mathrm{T}}(𝐫_1)C\gamma _5\overline{d}^\mathrm{b}(𝐫_1)\right).`$ (74) Integrating over both delta-functions and then using the normalization (13) we obtain $`{\displaystyle d^3𝐫_1\mathrm{e}^{i𝐤_1𝐫_1}d^3𝐫_2\mathrm{e}^{i𝐤_2𝐫_2}}`$ $`\times 0\|[\widehat{\psi }_p^{\alpha _2}(𝐫_2,0),[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}]_+\|0`$ $`=2(\gamma _0)_{\beta \alpha _1}(\gamma _0)_{\alpha \alpha _2},`$ (75) and with (72) $`p(k_2,\sigma _2)\|\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i}\|p(k_1,\sigma _1)=2\overline{u}_p^\alpha (k_2,\sigma _2)u_p^\beta (k_1,\sigma _1).`$ We note an analog relation for the d quark $`p(k_2,\sigma _2)\|\widehat{\overline{d}}_\alpha ^\mathrm{i}\widehat{d}_\beta ^\mathrm{i}\|p(k_1,\sigma _1)`$ $`=`$ $`1\overline{u}_p^\alpha (k_2,\sigma _2)u_p^\beta (k_1,\sigma _1),`$ while for the neutron we have $`n(k_2,\sigma _2)\|\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i}\|n(k_1,\sigma _1)`$ $`=`$ $`1\overline{u}_n^\alpha (k_2,\sigma _2)u_n^\beta (k_1,\sigma _1),`$ (78) $`n(k_2,\sigma _2)\|\widehat{\overline{d}}_\alpha ^\mathrm{i}\widehat{d}_\beta ^\mathrm{i}\|n(k_1,\sigma _1)`$ $`=`$ $`2\overline{u}_n^\alpha (k_2,\sigma _2)u_n^\beta (k_1,\sigma _1).`$ Using relations like $`[\widehat{A},\widehat{B}\widehat{C}]_+=[\widehat{A},\widehat{B}]_+\widehat{C}\widehat{B}[\widehat{A},\widehat{C}]_{}`$ analog equations for the four-quark condensates can be obtained. Two illustrative examples are given for the flavor-unmixed four-quark condensate for the proton $`{\displaystyle d^3𝐫_1\mathrm{e}^{i𝐤_1𝐫_1}d^3𝐫_2\mathrm{e}^{i𝐤_2𝐫_2}}`$ $`\times 0\|[\widehat{\psi }_p^{\alpha _2}(𝐫_2,0),[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{i}\widehat{\overline{u}}_\gamma ^\mathrm{j}\widehat{u}_\delta ^\mathrm{j},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}]_+\|0`$ $`={\displaystyle \frac{1}{6}}\widehat{\overline{u}}\widehat{u}_0(3(\gamma _0)_{\delta \alpha _1}(\gamma _0)_{\alpha _2\gamma }\delta _{\alpha \beta }`$ $`(\gamma _0)_{\delta \alpha _1}(\gamma _0)_{\alpha _2\alpha }\delta _{\beta \gamma }+3(\gamma _0)_{\beta \alpha _1}(\gamma _0)_{\alpha _2\alpha }\delta _{\gamma \delta }`$ $`(\gamma _0)_{\beta \alpha _1}(\gamma _0)_{\alpha _2\gamma }\delta _{\alpha \delta }),`$ (80) where the normalization (13) and $`\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{j}_0=\frac{1}{12}\delta ^{\mathrm{i}\mathrm{j}}\delta _{\alpha \beta }\widehat{\overline{u}}\widehat{u}_0`$ sumrule has been used. For four-quark condensates with Gell-Mann matrices involved we find $`{\displaystyle d^3𝐫_1\mathrm{e}^{i𝐤_1𝐫_1}d^3𝐫_2\mathrm{e}^{i𝐤_2𝐫_2}0\|[\widehat{\psi }_p^{\alpha _2}(𝐫_2,0),[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{j}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\delta ^\mathrm{l},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}]_+\|0\left(\lambda ^a\right)^{\mathrm{i}\mathrm{j}}\left(\lambda ^a\right)^{\mathrm{k}\mathrm{l}}}`$ $`=\|A_p\|^2ϵ^{\mathrm{a}\mathrm{b}\mathrm{c}}ϵ^{\mathrm{a}^{}\mathrm{b}^{}\mathrm{c}^{}}\left(u^{\mathrm{a}^{}\mathrm{T}}C\gamma _5d^\mathrm{b}^{}\right)\left(\overline{u}^{\mathrm{a}\mathrm{T}}C\gamma _5\overline{d}^\mathrm{b}\right)\left(2\delta ^{\mathrm{i}\mathrm{l}}\delta ^{\mathrm{k}\mathrm{j}}{\displaystyle \frac{2}{3}}\delta ^{\mathrm{i}\mathrm{j}}\delta ^{\mathrm{k}\mathrm{l}}\right)`$ $`\times ((\gamma _0)_{\delta \alpha _1}(\gamma _0)_{\alpha _2\gamma }\delta _{\alpha \beta }\delta ^{\mathrm{c}\mathrm{l}}\delta ^{\mathrm{c}^{}\mathrm{k}}\delta ^{\mathrm{i}\mathrm{j}}(\gamma _0)_{\delta \alpha _1}(\gamma _0)_{\alpha _2\alpha }\delta _{\beta \gamma }\delta ^{\mathrm{c}\mathrm{l}}\delta ^{\mathrm{i}\mathrm{c}^{}}\delta ^{\mathrm{j}\mathrm{k}}`$ $`+(\gamma _0)_{\beta \alpha _1}(\gamma _0)_{\alpha _2\alpha }\delta _{\gamma \delta }\delta ^{\mathrm{j}\mathrm{c}}\delta ^{\mathrm{i}\mathrm{c}^{}}\delta ^{\mathrm{k}\mathrm{l}}(\gamma _0)_{\beta \alpha _1}(\gamma _0)_{\alpha _2\gamma }\delta _{\alpha \delta }\delta ^{\mathrm{j}\mathrm{c}}\delta ^{\mathrm{c}^{}\mathrm{k}}\delta ^{\mathrm{i}\mathrm{l}})`$ $`={\displaystyle \frac{8}{9}}\left[(\gamma _0)_{\delta \alpha _1}(\gamma _0)_{\alpha _2\alpha }\delta _{\beta \gamma }+(\gamma _0)_{\beta \alpha _1}(\gamma _0)_{\alpha _2\gamma }\delta _{\alpha \delta }\right]\widehat{\overline{u}}\widehat{u}_0.`$ (81) By using the same technique the general result for the six-quark condensates is obtained as $`{\displaystyle d^3𝐫_1\mathrm{e}^{i𝐤_1𝐫_1}d^3𝐫_2\mathrm{e}^{i𝐤_2𝐫_2}0\|[\widehat{\psi }_p^{\alpha _2}(𝐫_2,0),[\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{j}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\delta ^\mathrm{l}\widehat{\overline{u}}_ϵ^\mathrm{m}\widehat{u}_\zeta ^\mathrm{n},\widehat{\overline{\psi }}_p^{\alpha _1}(𝐫_1,0)]_{}]_+\|0}`$ $`=\|A_p\|^2ϵ^{\mathrm{a}\mathrm{b}\mathrm{c}}ϵ^{\mathrm{a}^{}\mathrm{b}^{}\mathrm{c}^{}}\left(u^{\mathrm{a}^{}\mathrm{T}}C\gamma _5d^\mathrm{b}^{}\right)\left(\overline{u}^{\mathrm{a}\mathrm{T}}C\gamma _5\overline{d}^\mathrm{b}\right)(\widehat{u}_\beta ^\mathrm{j}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\delta ^\mathrm{l}\widehat{\overline{u}}_ϵ^\mathrm{m}_0(\gamma _0)_{\alpha _1\zeta }(\gamma _0)_{\alpha _2\alpha }\delta ^{\mathrm{c}\mathrm{n}}\delta ^{\mathrm{c}^{}\mathrm{i}}`$ $`+\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{j}\widehat{u}_\delta ^\mathrm{l}\widehat{\overline{u}}_ϵ^\mathrm{m}_0(\gamma _0)_{\alpha _1\zeta }(\gamma _0)_{\alpha _2\gamma }\delta ^{\mathrm{c}\mathrm{n}}\delta ^{\mathrm{c}^{}\mathrm{k}}+\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{j}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\delta ^\mathrm{l}_0(\gamma _0)_{\alpha _1\zeta }(\gamma _0)_{\alpha _2ϵ}\delta ^{\mathrm{c}\mathrm{n}}\delta ^{\mathrm{c}^{}\mathrm{m}}`$ $`+\widehat{u}_\beta ^\mathrm{j}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{\overline{u}}_ϵ^\mathrm{m}\widehat{u}_\zeta ^\mathrm{n}_0(\gamma _0)_{\alpha _1\delta }(\gamma _0)_{\alpha _2\alpha }\delta ^{\mathrm{c}\mathrm{l}}\delta ^{\mathrm{c}^{}\mathrm{i}}+\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{j}\widehat{\overline{u}}_ϵ^\mathrm{m}\widehat{u}_\zeta ^\mathrm{n}_0(\gamma _0)_{\alpha _1\delta }(\gamma _0)_{\alpha _2\gamma }\delta ^{\mathrm{c}\mathrm{l}}\delta ^{\mathrm{c}^{}\mathrm{k}}`$ $`\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\beta ^\mathrm{j}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\zeta ^\mathrm{n}_0(\gamma _0)_{\alpha _1\delta }(\gamma _0)_{\alpha _2ϵ}\delta ^{\mathrm{c}\mathrm{l}}\delta ^{\mathrm{c}^{}\mathrm{m}}+\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\delta ^\mathrm{l}\widehat{\overline{u}}_ϵ^\mathrm{m}\widehat{u}_\zeta ^\mathrm{n}_0(\gamma _0)_{\alpha _1\beta }(\gamma _0)_{\alpha _2\alpha }\delta ^{\mathrm{c}\mathrm{j}}\delta ^{\mathrm{c}^{}\mathrm{i}}`$ $`\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{u}_\delta ^\mathrm{l}\widehat{\overline{u}}_ϵ^\mathrm{m}\widehat{u}_\zeta ^\mathrm{n}_0(\gamma _0)_{\alpha _1\beta }(\gamma _0)_{\alpha _2\gamma }\delta ^{\mathrm{c}\mathrm{j}}\delta ^{\mathrm{c}^{}\mathrm{k}}\widehat{\overline{u}}_\alpha ^\mathrm{i}\widehat{\overline{u}}_\gamma ^\mathrm{k}\widehat{u}_\delta ^\mathrm{l}\widehat{u}_\zeta ^\mathrm{n}_0(\gamma _0)_{\alpha _1\beta }(\gamma _0)_{\alpha _2ϵ}\delta ^{\mathrm{c}\mathrm{j}}\delta ^{\mathrm{c}^{}\mathrm{m}}).`$ (82) Note that for applying the normalizations (13) and (14) one needs a term $`\delta ^{\mathrm{c}\mathrm{c}^{}}`$, which naturally arises when evaluating a specific matrix element under consideration.
warning/0506/hep-ph0506116.html	ar5iv	text	# Electron- and neutrino-nucleus scattering in the impulse approximation regime ## I Introduction The field of neutrino physics is rapidly developing after atmospheric neutrino oscillations and solar neutrino oscillations have been established Kajita ; Solar ; KamLAND . Recently, the SK Collaboration has found evidence of the oscillatory signature in atmospheric neutrinos, improving the determination of $`\mathrm{\Delta }m^2`$ SKATM , and K2K experiment has confirmed the oscillations of atmospheric neutrinos at 99.995% CL K2K ; Nakaya . These neutrino experiments measure energy and angle of muons produced in neutrino-nucleus interactions and reconstruct the incident neutrino energy, which determines the neutrino oscillations. K2K took data in the $`E_\nu =0.53`$ GeV region, and the recent L/E analysis of the SK atmospheric neutrinos is mainly based on the dataset extending from 0.5 to 25 GeV. JPARC and NuMI neutrino experiments JPARC ; Nova propose to measure $`\nu _\mu \nu _e`$ oscillations and determine $`\mathrm{\Delta }m^2`$ with $`1\%`$ accuracy and $`\mathrm{sin}^22\theta _{13}`$ above 0.006, using a narrow-band neutrino beam at $`E_\nu =0.8`$ GeV (JPARC) and 2.0 GeV (NuMI, off-axis). In view of these developments, it is vital that theoretical calculations of cross sections and spectra achieve an accuracy comparable to the experimental one, which in turn requires that the nuclear response to weak interactions be under control at a quantitative level. At $`E_\nu `$=3 GeV or less, quasi-elastic scattering and quasi-free $`\mathrm{\Delta }`$ production are the dominant neutrino-nucleus processes. However, reactions in this energy regime are associated with a wide range of momentum transfer, thus involving different aspects of nuclear structure. Four decades of electron-nucleus scattering experiments have unequivocally shown that the mean-field approximation, underlying the nuclear shell model, does not provide a fully quantitative account of the data (see, e.g., Ref. omar\_nuint04 and references therein). When the momentum transfer involved is large, dynamical nucleon-nucleon (NN) correlations are known to be important, and a description of nuclear structure beyond the mean-field picture is needed. On the other hand, neutrino-nucleus reactions also occur, in fact rather appreciably, with a small momentum transfer. Comparison between the data at $`Q^2<0.2\mathrm{GeV}^2`$ and the predictions of the Fermi gas (FG) model Smith , showing a sizable deficit of events Nakaya ; Ishida , suggests that a more realistic description of both nuclear properties and the reaction mechanism is indeed required. In this paper we discuss the extension of the many-body theory of electron-nucleus scattering (see, e.g., Ref. vijay\_nuint01 and references therein) to the case of neutrino-induced reactions. We focus on the energy range $`0.71.2`$ GeV and analyze inclusive scattering of both electrons and neutrinos off oxygen, the main target nucleus in SK, K2K and other experiments. The quasi-elastic and quasi-free $`\mathrm{\Delta }`$ production cross sections obtained from the FG model Smith ; Seki are compared to the results of the many-body approach developed in Ref. gofsix , extensively used to analyze electron scattering data at beam energy up to few GeV gofsix ; bp ; bffs . Preliminary versions of the materials in this paper have appeared in the Proceedings of NuInt04 omar\_nuint04 ; seki\_nuint04 ; farina\_nuint04 . Section II is devoted to a summary of the formalism employed to calculate the electron-nucleus cross section at high momentum transfer, as well as to the discussion of the main ingredients entering its definition: the nuclear spectral function, the elementary cross section describing electron scattering off a bound nucleon and the folding function embodying the main effects of final state interactions. In Section III the results of our approach are compared to inclusive electron scattering data at $`0.2\stackrel{<}{_{}}Q^2\stackrel{<}{_{}}0.6`$ GeV<sup>2</sup>, while in Section IV we outline the extension of the formalism to the case of charged current neutrino-nucleus scettering. Finally, our conclusions are stated in Section V. ## II Many-body theory of the electroweak nuclear response ### II.1 Electron-nucleus cross section The differential cross section of the process $$e+Ae^{}+X,$$ (1) in which an electron carrying initial four-momentum $`k(E_e,𝐤)`$ scatters off a nuclear target to a state of four-momentum $`k^{}(E_e^{},𝐤^{})`$, the target final state being undetected, can be written in Born approximation as (see, e.g., Ref. IZ ) $$\frac{d^2\sigma }{d\mathrm{\Omega }_e^{}dE_e^{}}=\frac{\alpha ^2}{Q^4}\frac{E_e^{}}{E_e}L_{\mu \nu }W^{\mu \nu },$$ (2) where $`\alpha `$ is the fine structure constant and $`Q^2=q^2=𝐪^2\nu ^2`$, $`q=kk^{}(\nu ,𝐪)`$ being the four momentum transfer. The leptonic tensor, that can be written, neglecting the lepton mass, as $$L_{\mu \nu }=2\left[k_\mu k_\nu ^{}+k_\nu k_\mu ^{}g_{\mu \nu }(kk^{})\right],$$ (3) is completely determined by electron kinematics, whereas the nuclear tensor $`W^{\mu \nu }`$ contains all the information on target structure. Its definition involves the initial and final hadronic states $`\|0`$ and $`\|X`$, carrying four-momenta $`p_0`$ and $`p_X`$, respectively, as well as the nuclear electromagnetic current operator $`J^\mu `$: $$W^{\mu \nu }=\underset{X}{}0\|J^\mu \|XX\|J^\nu \|0\delta ^{(4)}(p_0+qp_X),$$ (4) where the sum includes all hadronic final states. Calculations of $`W^{\mu \nu }`$ at moderate momentum transfers $`(\|𝐪\|<0.5\mathrm{GeV})`$ can be carried out within nuclear many-body theory (NMBT), using nonrelativistic wave functions to describe the initial and final states and expanding the current operator in powers of $`\|𝐪\|/m`$, $`m`$ being the nucleon mass (see, e.g., Ref. rocco ). On the other hand, at higher values of $`\|𝐪\|`$, corresponding to beam energies larger than $`1`$ GeV, the description of the final states $`\|X`$ in terms of nonrelativistic nucleons is no longer possible. Calculations of $`W^{\mu \nu }`$ in this regime require a set of simplifying assumptions, allowing one to take into account the relativistic motion of final state particles carrying momenta $`𝐪`$ as well as the occurrence of inelastic processes, leading to the appearance of hadrons other than protons and neutrons. ### II.2 The impulse approximation The main assumptions underlying the impulse approximation (IA) scheme are that i) as the spatial resolution of a probe delivering momentum $`𝐪`$ is $`1/\|𝐪\|`$, at large enough $`\|𝐪\|`$ the target nucleus is seen by the probe as a collection of individual nucleons and ii) the particles produced at the interaction vertex and the recoiling ($`\mathrm{A}1`$)-nucleon system evolve indipendently of one another, which amounts to neglecting both statistical correlations due to Pauli blocking and dynamical Final State Interactions (FSI), i.e. rescattering processes driven by strong interactions. In the IA regime the scattering process off a nuclear target reduces to the incoherent sum of elementary processes involving only one nucleon, as schematically illustrated in Fig. 1. Within this picture, the nuclear current can be written as a sum of one-body currents $$J^\mu \underset{i}{}j_i^\mu ,$$ (5) while the final state reduces to the direct product of the hadronic state produced at the electromagnetic vertex, carrying momentum $`𝐩_x`$ and the $`(A1)`$-nucleon residual system, carrying momentum $`𝐩_{}=𝐪𝐩_x`$ (for simplicity, we omit spin indices) $$\|X\|x,𝐩_x\|,𝐩_{}.$$ (6) Using Eq. (6) we can rewrite the sum in Eq. (4) replacing $`{\displaystyle \underset{X}{}}\|XX\|`$ $``$ $`{\displaystyle \underset{x}{}}{\displaystyle d^3p_x\|x,𝐩_x𝐩_x,x\|}`$ (7) $`\times `$ $`{\displaystyle \underset{}{}}d^3p_{}\|,𝐩_{}𝐩_{},\|.`$ Substitution of Eqs. (5)-(7) into Eq. (4) and insertion of a complete set of free nucleon states, satisfying $$d^3p\|\mathrm{N},𝐩𝐩,\mathrm{N}\|=I,$$ (8) results in the factorization of the current matrix element $`0\|J^\mu \|X`$ $`=`$ $`{\displaystyle \frac{m}{\sqrt{𝐩_{}^2+m^2}}}0\|,𝐩_{};\mathrm{N},𝐩_{}`$ (9) $`\times `$ $`{\displaystyle \underset{i}{}}𝐩_{},N\|j_i^\mu \|x,𝐩_x,`$ leading to $`W^{\mu \nu }`$ $`=`$ $`{\displaystyle \underset{x,}{}}{\displaystyle d^3p_{}d^3p_x\|0\|,𝐩_{};\mathrm{N},𝐩_{}\|^2\frac{m}{E_𝐩_{}}}`$ $`\times `$ $`{\displaystyle \underset{i}{}}𝐩_{},\mathrm{N}\|j_i^\mu \|x,𝐩_x𝐩_x,x\|j_i^\nu \|\mathrm{N},𝐩_{}`$ $`\times `$ $`\delta ^{(3)}(𝐪𝐩_{}𝐩_x)\delta (\nu +E_0E_{}E_x),`$ where $`E_𝐩_{}=\sqrt{\|𝐩_{}\|^2+m^2}`$. Finally, using the identity $`\delta (\nu +E_0E_{}E_x)`$ $`=`$ $`{\displaystyle 𝑑E\delta (Em+E_0E_{})}`$ (11) $`\times `$ $`\delta (\nu E+mE_x),`$ and defining the target spectral function as foot1 $`P(𝐩,E)`$ $`=`$ $`{\displaystyle \underset{}{}}\|0\|,𝐩;\mathrm{N},𝐩\|^2`$ (12) $`\times \delta (Em+E_0E_{}),`$ we can rewrite Eq. (4) in the form $`W^{\mu \nu }(𝐪,\nu )`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle d^3p𝑑Ew_i^{\mu \nu }(\stackrel{~}{q})}`$ (13) $`\times \left({\displaystyle \frac{m}{E_𝐩}}\right)P(𝐩,E),`$ with $`E_𝐩=\sqrt{\|𝐩^2\|+m^2}`$ and $`w_i^{\mu \nu }`$ $`=`$ $`{\displaystyle \underset{x}{}}𝐩,\mathrm{N}\|j_i^\mu \|x,𝐩+𝐪𝐩+𝐪,x\|j_i^\nu \|\mathrm{N},𝐩`$ (14) $`\times `$ $`\delta (\stackrel{~}{\nu }+\sqrt{𝐩^2+m^2}E_x).`$ The quantity defined in the above equation is the tensor describing electromagnetic interactions of the $`i`$-th nucleon in free space. Hence, Eq. (14) shows that in the IA scheme the effect of nuclear binding of the struck nucleon is accounted for by the replacement $$q(\nu ,𝐪)\stackrel{~}{q}(\stackrel{~}{\nu },𝐪),$$ (15) with (see Eqs. (II.2) and (12)) $`\stackrel{~}{\nu }`$ $`=`$ $`E_x\sqrt{𝐩^2+m^2}`$ (16) $`=`$ $`\nu +E_0E_{}\sqrt{𝐩^2+m^2}`$ $`=`$ $`\nu E+m\sqrt{𝐩^2+m^2},`$ in the argument of $`w_i^{\mu \nu }`$. This procedure essentially amounts to assuming that: i) a fraction $`\delta \nu `$ of the energy transfer goes into excitation energy of the spectator system and ii) the elementary scattering process can be described as if it took place in free space with energy transfer $`\stackrel{~}{\nu }=\nu \delta \nu `$. This interpretation emerges most naturally in the $`\|𝐩\|m`$ limit, in which Eq. (16) yields $`\delta \nu =E`$. Collecting together all the above results we can finally rewrite the doubly differential nuclear cross section in the form $`{\displaystyle \frac{d\sigma _{IA}}{d\mathrm{\Omega }_e^{}dE_e^{}}}={\displaystyle }d^3pdEP(𝐩,E)[Z{\displaystyle \frac{d\sigma _{ep}}{d\mathrm{\Omega }_e^{}dE_e^{}}}`$ $`+(AZ){\displaystyle \frac{d\sigma _{en}}{d\mathrm{\Omega }_e^{}dE_e^{}}}]\delta (\nu E+mE_x),`$ (17) where $`d\sigma _{eN}/d\mathrm{\Omega }_e^{}dE_e^{}`$ ($`Nn,p`$ denotes a proton or a neutron) is the cross section describing the elementary scattering process $$e(k)+N(p)e^{}(k^{})+x(p+\stackrel{~}{q}),$$ (18) given by $$\frac{d\sigma _{eN}}{d\mathrm{\Omega }_e^{}dE_e^{}}=\frac{\alpha ^2}{Q^4}\frac{E_e^{}}{E_e}\frac{m}{E_𝐩}L_{\mu \nu }w_N^{\mu \nu },$$ (19) stripped of both the flux factor and the energy conserving $`\delta `$-function. ### II.3 The nuclear spectral function In NMBT the nucleus is seen as a system of A nucleons whose dynamics are described by the nonrelativistic hamiltonian $$H_A=\underset{i=1}{\overset{A}{}}\frac{𝐩_i^2}{2m}+\underset{j>i=1}{\overset{A}{}}v_{ij}+\underset{k>j>i=1}{\overset{A}{}}V_{ijk},$$ (20) where $`𝐩_i`$ is the momentum of the $`i`$-th nucleon, while $`v_{ij}`$ and $`V_{ijk}`$ are two- and three-nucleon interaction potentials, respectively. The two-nucleon potential, that reduces to the Yukawa one-pion-exchange potential at large internucleon distance, is obtained from an accurate fit to the available data on the two-nucleon system, i.e. deuteron properties and $``$ 4000 NN scattering phase shifts at energies up to the pion production threshold WSS . The additional three-body term $`V_{ijk}`$ has to be included in order to account for the binding energies of the three-nucleon bound states PPCPW . The many-body Schrödinger equation associated with the Hamiltonian of Eq. (20) can be solved exactly, using stochastic methods, for nuclei with mass number $`A10`$. The energies of the ground and low-lying excited states are in excellent agreement with the experimental data WP . Accurate calculations can also be carried out for uniform nucleon matter, exploiting translational invariace and using either a variational approach based on cluster expansion and chain summation techniques AP , or G-matrix perturbation theory BGLS2000 . Nonrelativistic NMBT provides a fully consistent computational framework that has been employed to obtain the spectral functions of the few-nucleon systems, having A$`=3`$ dieperink ; cps ; sauer and 4 ciofi4 ; morita ; bp , as well as of nuclear matter, i.e. in the limit A $`\mathrm{}`$ with Z=A/2 bff ; pkebbg . Calculations based on G-matrix perturbation theory have also been carried out for oxygen geurts16 ; polls16 . The spectral functions of different nuclei, ranging from Carbon to Gold, have been modeled using the Local Density Approximation (LDA) bffs , in which the experimental information obtained from nucleon knock-out measurements is combined with the results of theoretical calculations of the nuclear matter $`P(𝐩,E)`$ at different densities bffs . Nucleon removal from shell model states has been extensively studied by coincidence $`(e,e^{}p)`$ experiments (see, e.g., Ref. book ). The corresponding measured spectral function is usually written in the factorized form $$P_{MF}(𝐩,E)=\underset{n}{}Z_n\|\varphi _n(𝐩)\|^2F_n(EE_n),$$ (21) where $`\varphi _n(𝐩)`$ is the momentum-space wave function of the single particle shell mode state $`n`$, whose energy width is described by the function $`F_n(EE_n)`$. The normalization of the $`n`$-th state is given by the so called spectroscopic factor $`Z_n<1`$, and the sum in Eq. (21) is extended to all occupied states of the Fermi sea. Hence, $`P_{MF}(𝐩,E)`$ vanishes at $`\|𝐩\|`$ larger than the Fermi momentum $`p_F250`$ MeV. Note that in absence of NN correlations $`F_n(EE_n)`$ shrinks to a $`\delta `$-function, $`Z_n1`$ and Eq. (21) can be identified with the full spectral function. Strong dynamical NN correlations give rise to virtual scattering processes leading to the excitation of the participating nucleons to states of energy larger than the Fermi energy, thus depleting the shell model states within the Fermi sea. As a consequence, the spectral function associated with nucleons belonging to correlated pairs extends to the region of $`\|𝐩\|p_F`$ and $`Ee_F`$, where $`e_F`$ denotes the Fermi energy, typically $`\stackrel{<}{_{}}30`$ MeV. The correlation contribution to $`P(𝐩,E)`$ of uniform nuclear matter has been calculated by Benhar et al for a wide range of density values bffs . Within the LDA scheme, the results of Ref. bffs can be used to obtain the corresponding quantity for a finite nucleus of mass number $`A`$ from $$P_{corr}(𝐩,E)=d^3r\rho _A(𝐫)P_{corr}^{NM}(𝐩,E;\rho =\rho _A(𝐫)),$$ (22) where $`\rho _A(𝐫)`$ is the nuclear density distribution and $`P_{corr}^{NM}(𝐩,E;\rho )`$ is the correlation component of the spectral function of uniform nuclear matter at density $`\rho `$. Finally, the full LDA nuclear spectral function can be written $$P_{LDA}(𝐩,E)=P_{MF}(𝐩,E)+P_{corr}(𝐩,E),$$ (23) the spectroscopic factors $`Z_n`$ of Eq. (21) being constrained by the normalization requirement $$d^3p𝑑EP_{LDA}(𝐩,E)=1.$$ (24) The LDA spectral function of $`{}_{}{}^{16}O`$ obtained combining the nuclear matter results of Ref. bffs and the Saclay $`(e,e^{}p)`$ data eep16O is shown in Fig. 2. The shell model contribution $`P_{MF}(𝐩,E)`$ accounts for $``$ 80 % of its normalization, whereas the remaining $``$ 20 % of the strength, accounted for by $`P_{corr}(𝐩,E)`$, is located at high momentum ($`\|𝐩\|p_F`$) and large removal energy ($`Ee_F`$). It has to be emphasized that large $`E`$ and large $`𝐩`$ are strongly correlated. For example, $``$ 50 % of the strength at $`\|𝐩\|`$ = 320 MeV is located at $`E>`$ 80 MeV. The LDA scheme rests on the premise that short range nuclear dynamics is unaffected by surface and shell effects. The validity of this assumption is confirmed by theoretical calculations of the nucleon momentum distribution, defined as $`n(𝐩)`$ $`=`$ $`{\displaystyle 𝑑EP(𝐩,E)}`$ (25) $`=`$ $`0\|a_𝐩^{}a_𝐩\|0,`$ (26) where $`a_𝐩^{}`$ and $`a_𝐩`$ denote the creation and annihilation operators of a nucleon of momentum $`𝐩`$. The results clearly show that for A$`4`$ the quantity $`n(𝐩)/A`$ becomes nearly independent of $`A`$ in the region of large $`\|𝐩\|`$ ($`\stackrel{>}{_{}}300`$ MeV), where NN correlations dominate (see, e.g., Ref. rmp ). In Fig. 3 the nucleon momentum distribution of <sup>16</sup>O, obtained from Eq. (25) using the LDA spectral function of Fig. 2, is compared to the one resulting from a Monte Carlo calculation steve , carried out using the definition of Eq. (26) and a highly realistic many-body wave function 16Owf . For reference, the FG model momemtum distribution corresponding to Fermi momentum $`p_F`$ = 221 MeV is also shown by the dashed line. It clearly appears that the $`n(𝐩)`$ obtained from the spectral function is close to that of Ref.steve , while the FG distribution exhibits a completely different behaviour. A direct measuremet of the correlation component of the spectral function of $`{}_{}{}^{12}C`$, obtained measuring the $`(e,e^{}p)`$ cross section at missing momentum and energy up to $``$ 800 MeV and $`200`$ MeV, respectively, has been recently carried out at Jefferson Lab by the E97-006 Collaboration E97-006 . The data resulting from the preliminary analysis appear to be consistent with the theoretical predictions based on LDA. ### II.4 Final state interactions The occurrence of strong FSI in electron nucleus scattering has long been experimentally established. The results of a number of $`(e,e^{}p)`$ measurements covering the kinematical domain corresponding to $`0.5\stackrel{<}{_{}}Q^2\stackrel{<}{_{}}8.0`$ GeV<sup>2</sup> garino92 ; oneill95 ; abbott98 ; garrow02 , clearly show that the flux of outgoing protons is strongly suppressed, with respect to the IA predictions. The observed attenuation ranges from 20-40 % in Carbon to 50-70 % in Gold. The inclusive $`(e,e^{})`$ cross section, being only sensitive to rescattering processes taking place within a distance $`1/\|𝐪\|`$ of the electromagnetic vertex, is obviously much less sensitive to FSI than the coincidence $`(e,e^{}p)`$ cross section. The latter is in fact affected by rescatterings occurring over the distance $`R_A`$, $`R_A`$ being the nuclear radius, travelled by the struck particle on its way out of the target. However, FSI effects become appreciable, indeed dominant, in the low $`\nu `$ region, where the inclusive cross section is most sensitive to the high momentum and high removal energy tails of the nuclear spectral function. In quasi-elastic inclusive processes FSI produce two effects: i) an energy shift of the cross section, due to the fact that the struck nucleon feels the mean field generated by the spectator particles and ii) a redistribution of the strength, leading to the quenching of the peak and the enhancement of the tails, to be ascribed to the occurrence of NN rescattering processes that couple the one particle-one hole final state to more complicated $`n`$ particle-$`n`$ holes configurations. Early attempts to include FSI effects were based on the optical potential model horikawa80 . However, while providing a computationally practical scheme to account for the loss of flux in the one-nucleon removal channel, this model relies on the mean field picture of the nucleus, and does not include the effect of dynamical NN correlations. A different approach, based on NMBT and a generalization of Glauber theory of high energy proton scattering glauber59 has been proposed by Benhar et al. gofsix in the early 90s. This treatment of FSI, generally referred to as Correlated Glauber Approximation (CGA) rests on the assumptions that i) the struck nucleon moves along a straight trajectory with constant velocity (eikonal approximation), and ii) the spectator nucleons are seen by the struck particle as a collection of fixed scattering centers (frozen approximation). Under the above assumptions the expectation value of the propagator of the struck nucleon in the target ground state can be written in the factorized form $$U_{𝐩+𝐪}(t)=U_{𝐩+𝐪}^0(t)\overline{U}_{𝐩+𝐪}^{FSI}(t),$$ (27) where $`U_{𝐩+𝐪}^0`$ is the free space propagator, while FSI effects are described by the quantity ($`R(𝐫_1,\mathrm{},𝐫_A)`$ specifies the target configuration) $$\overline{U}_{𝐩+𝐪}^{FSI}(t)=0\|U_{𝐩+𝐪}^{FSI}(R;t)\|0,$$ (28) with $$U_{𝐩+𝐪}^{FSI}(R;t)=\frac{1}{A}\underset{i=1}{\overset{A}{}}\mathrm{e}^{i_{ji}_0^t𝑑t^{}w_{𝐩+𝐪}(\|𝐫_{ij}+𝐯t^{}\|)}.$$ (29) In Eq. (29) $`𝐫_{ij}=𝐫_i𝐫_j`$ and $`w_{𝐩+𝐪}(\|𝐫\|)`$ is the coordinate-space NN scattering t-matrix at incident momentum $`𝐩+𝐪`$, usually parametrized in terms of total cross section, slope and real to imaginary part ratio. At large $`\|𝐪\|`$, $`𝐩+𝐪𝐪`$ and the eikonal propagator of Eq. (28) becomes a function of $`t`$ and the momentum transfer only. Note that $`U_𝐪^{FSI}(R;t)`$ is simply related to the nuclear transparency $`T_A`$, measured in coincidence $`(e,e^{}p)`$ experiments garino92 ; oneill95 ; abbott98 ; garrow02 , through $$T_A=\underset{t\mathrm{}}{lim}0\|\|U_𝐪^{FSI}(R;t)\|^2\|0.$$ (30) The results displayed in Fig. 4 daniela show that both the magnitude and the $`A`$\- and $`Q^2`$-dependence of the transparencies of Carbon, Iron and Gold obtained from the approach of Ref. gofsix and LDA are in good agreement with the experimental data. Note that at low $`Q^2`$ FSI lead to a $``$ 20 (40) % effect in Carbon (Iron). Neglecting this effect, i.e. setting $`T_A(Q^2)1`$, would be utterly incompatible with the data. From Eqs. (28) and (29) it follows that within the approach of Ref. gofsix the energy shift and the redistribution of the inclusive strength are driven by the real and the imaginary part of the NN scattering amplitude, respectively. At large $`𝐪`$ the imaginary part of $`w_𝐪`$, corresponding to the real part of $`U_𝐪^{FSI}`$, is dominant. Neglecting the contribution of the real part of $`w_𝐪`$ altogether, the CGA quasi-elastic inclusive cross section can be written as a convolution integral, involving the cross section evaluated within the IA, i.e. in absence of FSI, and a folding function embodying FSI effects: $$\frac{d\sigma }{d\mathrm{\Omega }_e^{}d\nu }=𝑑\nu ^{}f_𝐪(\nu \nu ^{})\left(\frac{d\sigma }{d\mathrm{\Omega }_e^{}d\nu ^{}}\right)_{IA},$$ (31) the folding function $`f_𝐪(\nu )`$ being defined as $$f_𝐪(\nu )=\delta (\nu )\sqrt{T_A}+\frac{dt}{2\pi }\mathrm{e}^{i\nu t}\left[U_𝐪^{FSI}(t)\sqrt{T_A}\right].$$ (32) The above equations clearly show that the strength of FSI is measured by both $`T_A`$ and the width of the folding function. In absence of FSI $`U_𝐪^{FSI}(R;t)1`$, implying in turn $`T_A=1`$ and $`f_𝐪(\nu )=\delta (\nu )`$. Dynamical NN correlations stronlgy affect the shape of the folding function of Eq. (32). Due to the strong repulsive core of the NN force, the joint probability of finding two nucleons at positions $`𝐫_i`$ and $`𝐫_j`$, driving the occurrence of rescattering processes in the final state, is strongly suppressed at $`\|𝐫_i𝐫_j\|\stackrel{<}{_{}}1`$ fm. As a consequence, inclusion of correlation effects within the framework of NMBT leads to a strong quenching of FSI effects, with respect to the predictions of the independent particle model. In principle, the real part of the NN scattering ampliutde can be explicitely included in Eq. (29) and treated on the same footing as the imaginary part. However, its effect turns out to be appreaciable only at $`t0`$, when the attenuation produced by the imaginary part is weak. The results of numerical calculations show that an approximate treatment based on the use of a time independent optical potential is indeed adequate to describe the energy shift produced by the real part of $`w_𝐪`$ bffs , whose size of $`10`$ MeV is to be compared to a typical electron energy loss of few hundreds MeV. ## III Comparison to electron scattering data We have employed the formalism described in the previuos Sections to compute the inclusive electron scattering cross section off oxygen at $`0.2\stackrel{<}{_{}}Q^2\stackrel{<}{_{}}0.6`$ GeV<sup>2</sup>. The IA cross section has been obtained using the LDA spectral function shown in Fig. 2 and the nucleon tensor defined by Eq. (14), that can be written as $`w_N^{\mu \nu }`$ $`=`$ $`w_1^N\left(g^{\mu \nu }+{\displaystyle \frac{\stackrel{~}{q}^\mu \stackrel{~}{q}^\nu }{\stackrel{~}{q}^2}}\right)`$ (33) $`+`$ $`{\displaystyle \frac{w_2^N}{m^2}}\left(p^\mu {\displaystyle \frac{(p\stackrel{~}{q})}{\stackrel{~}{q}^2}}q^\mu \right)\left(p^\nu {\displaystyle \frac{(p\stackrel{~}{q})}{\stackrel{~}{q}^2}}q^\nu \right),`$ where $`p(E_𝐩,𝐤)`$ and the off-shell four momentum transfer $`\stackrel{~}{q}`$ is defined by Eqs. (15) and (16). The two structure functions $`w_1^N`$ and $`w_2^N`$ are extracted from electron-proton and electron-deuteron scattering data. In the case of quasi-elastic scattering they are simply related to the electric and magnetic nucleon form factors, $`G_{E_N}`$ and $`G_{M_N}`$, through $$w_1^N=\frac{\stackrel{~}{q}^2}{4m^2}\delta \left(\stackrel{~}{\nu }+\frac{\stackrel{~}{q}^2}{2m}\right)G_{M_N}^2,$$ (34) $`w_2^N`$ $`=`$ $`{\displaystyle \frac{1}{1\stackrel{~}{q}^2/4m^2}}\delta \left(\stackrel{~}{\nu }+{\displaystyle \frac{\stackrel{~}{q}^2}{2m}}\right)`$ (35) $`\times \left(G_{E_N}^2{\displaystyle \frac{\stackrel{~}{q}^2}{4m^2}}G_{M_N}^2\right).`$ Numerical calculations have been carried out using the Höhler-Brash parameterization of the form factors Hohler76 ; Brash02 , resulting from a fit which includes the recent Jefferson Lab data Jones00 . In the kinematical region under discussion, inelastic processes, mainly quasi-free $`\mathrm{\Delta }`$ resonance production, are also known to play a role. To include these contributions in the calculation of the inclusive cross section, we have adopted the Bodek and Ritchie parametrization of the proton and neutron structure functions br , covering both the resonance and deep inelastic region. The folding functions describing the effect of NN rescattering in the final state have been computed from Eq. (32) with the eikonal propagator $`U_𝐪^{FSI}(R;t)`$ obtained using the parametrization of the NN scattering amplitude of Ref. oneillNN and the medium modified NN cross sections of Ref. papi . The integrations involved in Eq. (28) have been carried out using Monte Carlo configurations sampled from the probability distribution associated with the oxygen ground state wave function of Ref. 16Owf . The effect of the real part of the NN scattering amplitude has been approximated including in the energy conserving $`\delta `$-function of Eq. (17) the real part of the optical potential felt by a nucleon of momentum $`𝐩+𝐪`$ embedded in uniform nuclear matter at equilibrium density. In Figs. 5-8 the results of our calculations are compared to the data of Ref. LNF , corresponding to beam energies 700, 880, 1080 and 1200 MeV and electron scattering angle 32. For reference, the results of the FG model corresponding to Fermi momentum $`p_F=225`$ MeV and average removal energy $`ϵ=25`$ MeV are also shown. Overall, the approach described in the previous Sections, involving no adjustable parameters, provides a fairly accurate account of the measured cross sections in the region of the quasi-free peak. On the other hand, the FG model, while yielding a reasonable description at beam energies 1080 and 1200 MeV, largely overestimates the data at lower energies. The discrepancy at the top of the quasi-elastic peak turns out to be $``$ 25 % and $``$ 50 % at 880 and 700 MeV, respectively. The results of NMBT and FG model also turn out to be sizably different in the dip region, on the right hand side of the quasi-elastic peak, while the discrepancies become less pronounced at the $`\mathrm{\Delta }`$-production peak. However, it clearly appears that, independent of the employed approach and beam energy, theoretical results significantly underestimate the data at energy transfer larger than the pion production threshold. In view of the fact that the quasi-elastic peak is correctly reproduced (within an accuracy of $``$ 10 %), the failure of NMBT to reproduce the data at larger $`\omega `$ may be ascribed to deficiencies in the description of the elementary electron-nucleon cross section. In fact, as ilustrated in Fig. 9, the calculation of the IA cross section at the quasi-elastic and $`\mathrm{\Delta }`$ production peak involves integrations of $`P(𝐩,E)`$ extending over regions of the $`(𝐩,E)`$ plane almost exaclty overlapping one another. To gauge the uncertainty associated with the description of the nucleon structure functions $`w_1^N`$ and $`w_2^N`$, we have compared the electron-proton cross sections obtained from the model of Ref. br to the ones obtained from the $`H_2`$ model of Ref. thia and from a global fit including recent Jefferson Lab data christy . The results of Fig. 10 show that at $`E_e=1200`$ MeV and $`\theta =32^{}`$ the discrepancy between the different models is not large, being $``$ 15 % at the $`\mathrm{\Delta }`$ production peak. It has to be noticed, however, that the models of Refs. br ; thia ; christy have all been obtained fitting data taken at electron beam energies larger than 2 GeV, so that their use in the kinematical regime discussed in this work involves a degree of extrapolation. On the other hand, the results obtained using the approach described in this paper and the nucleon structure functions of Ref. br are in excellent agreement with the measured $`(e,e^{})`$ cross sections at beam energies of few GeV bffs . As an example, in Fig. 11 we show a comparison between the calculated $`{}_{}{}^{12}C(e,e^{})`$ cross section and the Jefferson Lab data of Ref. E89-008 at $`E_e=4`$ GeV and $`\theta =30^{}`$. The corresponding FG result is also displayed, for reference. Fig. 10 also shows the prediction of the Bodek and Ritchie fit for the neutron cross section, which turns out to be much smaller than the proton one. It is on account of this difference that we have chosen to adopt the fit of Ref. br , as it allows for a consistent inclusion of proton and neutron contributions, both resonant and nonresonsant, to the nuclear cross section. In this regard, it has to be pointed out that the nonresonant backround is not negligible. As illustrated in Fig. 12, for beam energy 1200 MeV and scattering angle 32 it provides $``$ 25 % of the cross section at energy transfer corresponding to the $`\mathrm{\Delta }`$ peak. The folding function described in Section II accounts for the FSI between a nucleon and the spectator system. For this reason, the results shown in Figs. 5-8 have been obtained folding with $`f_𝐪(\nu )`$ only the quasielastic component of the IA cross section. Particles other than protons and neutrons, that can be produced at the electromagnetic vertex, also have FSI, but they are more difficult to describe. However, the inelastic part of the IA cross section, being rather smooth, is unlikely to be strongly affected by FSI. To gauge the possible relevance of neglecting FSI in the inelastic channels we have computed the cross section at incident energy 880 MeV and scattering angle 32 folding the total IA result. Comparison with the results displayed in Fig. 6 shows a difference of $``$ 0.5 % at the top of the $`\mathrm{\Delta }`$-production peak. ## IV Charged current neutrino-nucleus cross section The Born approximation cross section of the weak charged current process $$\nu _{\mathrm{}}+A\mathrm{}^{}+X,$$ (36) can be written in the form (compare to Eq. (2)) $$\frac{d\sigma }{d\mathrm{\Omega }_{\mathrm{}}dE_{\mathrm{}}}=\frac{G^2}{32\pi ^2}\frac{\|𝐤^{}\|}{\|𝐤\|}L_{\mu \nu }W^{\mu \nu },$$ (37) where $`G=G_F\mathrm{cos}\theta _C`$, $`G_F`$ and $`\theta _C`$ being Fermi’s coupling constant and Cabibbo’s angle, $`E_{\mathrm{}}`$ is the energy of the final state lepton and $`𝐤`$ and $`𝐤^{}`$ are the neutrino and charged lepton momenta, respectively. Compared to the corresponding quantities appearing in Eq. (2), the tensors $`L_{\mu \nu }`$ and $`W^{\mu \nu }`$ include additional terms resulting from the presence of axial-vector components in the leptonic and hadronic currents (see, e.g., Ref. walecka ). Within the IA scheme, the cross section of Eq. (37) can be cast in a form similar to that obtained for the case of electron-nucleus scattering (see Eq. (17)). Hence, its calculation requires the nuclear spectral function and the tensor describing the weak charged current interaction of a free nucleon, $`w_N^{\mu \nu }`$. In the case of quasi-elastic scattering, neglecting the contribution associated with the pseudoscalar form factor $`F_P`$, the latter can be written in terms of the nucleon Dirac and Pauli form factors $`F_1`$ and $`F_2`$, related to the measured electric and magnetic form factors $`G_E`$ and $`G_M`$ through $$F_1=\frac{1}{1q^2/4m^2}\left(G_E\frac{q^2}{4m^2}G_M\right)$$ (38) $$F_2=\frac{1}{1q^2/4m^2}\left(G_MG_E\right),$$ (39) and the axial form factor $`F_A`$. Figure 13 shows the calculated cross section of the process $`{}_{}{}^{16}O(\nu _e,e)`$, corresponding to neutrino energy $`E_\nu =1`$ GeV and electron scattering angle $`\theta _e=30^{}`$, plotted as a function of the energy transfer $`\nu =E_\nu E_e`$. Numerical results have been obtained using the spectral function of Fig. 2 and the dipole parametrization for the form factors, with an axial mass of 1.03 GeV. Comparison between the solid and dashed lines shows that the inclusion of FSI results in a sizable redistribution of the IA strength, leading to a quenching of the quasi-elastic peak and to the enhancement of the tails. For reference, we also show the cross section predicted by the FG model with Fermi momentum $`p_F=225`$ MeV and average separation energy $`ϵ=25`$ MeV. Nuclear dynamics, neglected in the oversimplified picture in terms of noninteracting nucleons, clearly appears to play a relevant role. It has to be pointed out that the approach described in Section II, while including dynamical correlations in the final state, does not take into account statistical correlations, leading to Pauli blocking of the phase space available to the knocked-out nucleon. A rather crude prescription to estimate the effect of Pauli blocking amounts to modifying the spectral function through the replacement $$P(𝐩,E)P(𝐩,E)\theta (\|𝐩+𝐪\|\overline{p}_F)$$ (40) where $`\overline{p}_F`$ is the average nuclear Fermi momentum, defined as $$\overline{p}_F=d^3r\rho _A(𝐫)p_F(𝐫),$$ (41) with $`p_F(𝐫)=(3\pi ^2\rho _A(𝐫)/2)^{1/3}`$, $`\rho _A(𝐫)`$ being the nuclear density distribution. For oxygen, Eq. (41) yields $`\overline{p}_F=209`$ MeV. Note that, unlike the spectral function, the quantity defined in Eq. (40) does not describe intrinsic properties of the target only, as it depends explicitely on the momentum transfer. The effect of Pauli blocking is hardly visible in the differential cross section shown in Fig. 13, as the kinematical setup corresponds to $`Q^2>0.2`$ GeV<sup>2</sup> at the quasi-elastic peak. The same is true for the electron scattering cross sections discussed in the previous Section. On the other hand, this effect becomes very large at lower $`Q^2`$. Figure 14 shows the calculated differential cross section $`d\sigma /dQ^2`$ for neutrino energy $`E_\nu =1`$ GeV. The dashed and dot-dash lines correspond to the IA results with and without inclusion of Pauli blocking, respectively. It clearly appears that the effect of Fermi statistic in suppressing scattering shows up at $`Q^2<0.2`$ GeV<sup>2</sup> and becomes very large at lower $`Q^2`$. The results of the full calculation, in which dynamical FSI are also included, are displayed as a full line. The results of Fig. 14 suggest that Pauli blocking and FSI may explain the deficit of the measured cross section at low $`Q^2`$ with respect to the predictions of Monte Carlo simulations Ishida . Figure 15 shows the $`\nu _\mu `$-nucleus cross sections as a function of the scattered muon energy, by comparing the cross sections calculated by FG, and by the use of the spectral function with and without Pauli blocking. Figure 15 shows that FG yields a larger high-energy peak contribution than the other two. This is not due to the Pauli blocking, but due to the nuclear correlation effects in the spectral function: the muons tend to be scattered with a higher energy. This effect should show up in the forward angle cross section and may have a direct effect on neutrino oscillation measurements. ## V Conclusions We have employed an approach based on NMBT to compute the inclusive electron- and neutrino-nucleus scattering cross sections in the kinematical region corresponding to beam energy $``$ 1 GeV, relevant to many neutrino oscillation experiments. Our calculations have been carried out within the IA scheme, using realistic spectral functions obtained from $`(e,e^{}p)`$ data and theoretical calculations of uniform nuclear matter. In the region of the quasi-elastic peak, the results of our calculations account for the measured $`{}_{}{}^{16}O(e,e^{})`$ cross sections at beam energies between 700 MeV and 1200 MeV and scattering angle 32 with an accuracy better than 10 %. It must be emphasized that the ability to yield quantitative predictions over a wide range of beam energies is critical to the analysis of neutrino experiments, in which the energy of the incident neutrino is not known, and must be reconstructed from the kinematics of the outgoing lepton. In the region of quasi-free $`\mathrm{\Delta }`$ production theoretical predictions significantly underestimate the data. Assuming the validity of the IA scheme, this problem appears to be mainly ascribable to uncertainties in the description of the nucleon structure functions in this kinematical regime. The upcoming electron-nucleus scattetring data in the resonance region from the Jefferson Lab E04-001 experiment thia2 will help to shed light on this issue. At higher energies, i.e. in the region in which inelastic contributions largely dominate, the calculated cross sections are in close agreement with the data. Among the mechanisms not included in the IA picture, scattering processes in which the incoming lepton couples to meson exchange currents are not expected to produce large corrections to our results in the region of the quasi-elastic peak. Numerical studies of the transverse response of uniform nuclear matter, carried out within NMBT fabro , have shown that inclusion of two-body contributions to the nuclear electromagnetic current, arising from $`\pi `$ and $`\rho `$ meson exchange, leads to an enhancement that decreases as the momentum transfer increases, and never exceeds 10 % at $`Q^2<0.25`$ GeV<sup>2</sup>. On the other hand, the results of calculations of the tranverse response of the few-nucleon systems (for a review see Ref. rocco ) suggest that two-body current contributions may play a role in the dip region, at least for the lower values of the momentum transfer. The second mechanism not included in the IA, Pauli blocking, while not appreciably affecting the lepton energy loss spectra, produces a large effect on the $`Q^2`$ distributions at $`Q^2<0.2`$ GeV<sup>2</sup>, and must therefore be taken into account. In conclusion, NMBT provides a fully consistent and computationally viable scheme to calculate the electroweak nuclear response. Using the approach discussed in this paper may greatly contribute to decrease the systematic uncertainties associated with the analysis of neutrino oscillation experiments, as, unlike the FG model and other many-body approaches based on effective NN interactions (see, e.g., Ref. nieves and references therein), it is strongly constrained by NN data and involves no adjustable parameters. ###### Acknowledgements. This work is supported by the U. S. Department of Energy under grant DE-FG03-87ER40347 at CSUN and by the U. S. National Science Foundation under grant 0244899 at Caltech. One of the authors (OB) is deeply indebted to Vijay Pandharipande and Ingo Sick for a number of illuminating discussions on issues relevant to the subject of this work. Thanks are also due to Steven Pieper for providing Monte Carlo configurations sampled from the oxygen ground state wave function of Ref. 16Owf , as well as tables of the medium modified NN cross sections.
warning/0506/math-ph0506042.html	ar5iv	text	# Modulation of Camassa–Holm equation and reciprocal transformations ## 1 Introduction In 1993 R. Camassa and D. Holm proposed a new equation $$𝚞_t+3\mathrm{𝚞𝚞}_x=(𝚞_{xxt}+2𝚞_x𝚞_{xx}+\mathrm{𝚞𝚞}_{xxx})2\nu 𝚞_x,$$ (1) with $`\nu `$ a constant parameter, deriving it as the governing equation for waves in shallow water when surface tension is present. This involves an asymptotic expansion in small amplitude of the incompressible Euler equation for unidirectional motion under the influence of gravity that extends one order beyond the Korteweg-de Vries (KdV) equation. (1) is also an element in a class of equations introduced by A.Fokas and B.Fuchssteiner through the method of recursion operators in 1981. Equation (1) is strongly nonlinear, admits a bi-Hamiltonian structure , a Lax pair and it is formally integrable through the inverse scattering method . The bi-Hamiltonian structure of the CH equation can be described as follows $$m_t=P_1\frac{\delta H_2}{\delta m},m=𝚞𝚞_{xx},P_1=_x+_x^3,$$ with $$H_2=\frac{1}{2}(𝚞^3+\mathrm{𝚞𝚞}_x^2+2\nu 𝚞^2)𝑑x,$$ (2) or $$m_t=P_2\frac{\delta H_1}{\delta m},P_2=_xmm_x2\nu _x,$$ and $$H_1=\frac{1}{2}(𝚞^2+𝚞_x^2)𝑑x.$$ (3) The bi-Hamiltonian structure implies that the CH equation has an infinite number of conserved quantities. The functionals $`H_k`$, $`k,`$ defined by $$P_2\frac{\delta H_k}{\delta m}=P_1\frac{\delta H_{k+1}}{\delta m},k,$$ (4) are conserved quantities in involution with respect to the Poisson bracket determined either by $`P_1`$ or $`P_2`$. The Hamiltonian $`H_0=m𝑑x`$ is the Casimir of the first Poisson tensor $`P_1`$. For $`k>2`$ the Hamiltonian densities of $`H_k`$ are not local functions of $`𝚞`$ and its spatial derivatives. In the case $`\nu =0`$, R. Camassa and D. Holm proved the existence of solutions that are continuous but only piece-wise analytic (peakons). The CH equation possesses soliton solutions, periodic finite-gap solutions , real finite-gap solutions and, for $`\nu =0`$, multi-peakons . In particular, the algebro–geometric solutions of (CH) are described as Hamiltonian flows on nonlinear subvarieties (strata) of generalized Jacobians. This implies that the associated finite dimensional integrable systems may be described in the framework of integrable systems with deficiency whose algebraic–geometrical structure has much in common with the celebrated algebraically completely integrable systems introduced and thoroughly studied by M. Adler and P. van Moerbeke . In this work we derive the Whitham modulation equations for the CH flows. Whitham modulation equations for a nonlinear evolution system describe slow modulations of parameters over a family of periodic travelling wave solutions (or families of multi–phase solutions which are so far known to exist only for integrable systems). Contrary to the Korteweg de Vries case , it is an open problem to show that the Cauchy problem for CH with slowly varying initial data is described by the Whitham equations. Both for KdV and CH equations, this approximation is physically meaningful when the ratio between the water depth and the wavelength is very small . The Whitham equations are a system of hydrodynamic type equations and in the Riemann invariant coordinates take the form $$u_t^i+v^i(𝒖)u_x^i=0,i=1,\mathrm{},N,𝒖=(u^1,\mathrm{},u^N),$$ where we denote fast and slow variables with the same letters $`x`$ and $`t`$ and upper indeces denote controvariant vectors. The original evolution system is usually Lagrangian or Hamiltonian and this property is usually inherited by the equations of slow modulations. To average the original equations in the Lagrangian form, Whitham introduced the pseudo-phase method and then he constructed the corresponding Hamiltonian structure. For local Hamiltonian structures, B.A. Dubrovin and S.P. Novikov introduced a procedure for averaging local Poisson brackets and obtained the corresponding modulation equations. A third method to derive the Whitham modulation equations is the nonlinear analog of the WKB method . It can be proven that the three methods lead to the same equations for the case in which the original equation has a local Hamiltonian structure and local Hamiltonian densities (see and references therein). The Camassa-Holm equation does not possess a local Hamiltonian structure: indeed in the variable $`𝚞`$ the Hamiltonian operator is strongly nonlocal and the Hamiltonian densities of $`H_k`$ are non-local for $`k>2`$. In the variable $`m`$ the Hamiltonian densities of $`H_k`$ are nonlocal for $`k>0`$. Therefore the CH equation does not fit into the method of averaging local Hamiltonian structure nor even in the Maltsev-Novikov method of averaging weakly nonlocal Hamiltonian structures (an Hamiltonian structure is weakly nonlocal if it is polynomial in $`_x`$ and its higher derivatives and linear in $`_x^1`$) or in the Maltsev method of averaging weakly non-local symplectic form (inverse of the Hamiltonian operator). The latter method applies to Camassa-Holm only when averaging one–phase solutions . The CH equation can be written as a local Lagrangian system and we use the Whitham method (modulation equations in Lagrangian form) to derive the modulation equations for the one-phase periodic solution. The CH modulation equations for the Riemann invariants $`u^1<u^2<u^3`$, take the form $$_tu^i+C^i(𝒖)_xu^i=0,i=1,\mathrm{},3,$$ where $$\begin{array}{cc}& C^1(u^1,u^2,u^3)=u^1+u^2+u^3+2\nu +2\frac{(u^1+\nu )(u^1u^2)\mathrm{\Lambda }(K(s),\rho ,s)}{(u^2+\nu )[K(s)E(s)]}\hfill \\ & C^2(u^1,u^2,u^3)=u^1+u^2+u^3+2\nu +\frac{2(u^2u^1)\mathrm{\Lambda }(K(s),\rho ,s)}{K(s){\displaystyle \frac{(u^2+\nu )(u^3u^1)}{(u^1+\nu )(u^3u^2)}}E(s)}\hfill \\ & C^3(u^1,u^2,u^3)=u^1+u^2+u^3+2\nu +2\frac{(u^1+\nu )(u^3u^2)\mathrm{\Lambda }(K(s),\rho ,s)}{(u^2+\nu )E(s)}.\hfill \end{array}$$ In the above formulas $`K(s)`$ and $`E(s)`$ are the complete elliptic integrals of the first and second kind with modulus $`s^2={\displaystyle \frac{(u^2u^1)(u^3+\nu )}{(u^3u^1)(u^2+\nu )}}`$ and $`\mathrm{\Lambda }(K(s),\rho ,s)`$ is the complete elliptic integral of the third kind defined by $$\mathrm{\Lambda }(K(s),\rho ,s)=_0^{K(s)}\frac{dv}{1\rho ^2sn^2v},\rho ^2=\frac{u^2u^1}{u^2+\nu },$$ with $`sn`$ the Jacobi elliptic function. The equations are hyperbolic for $`\nu <u^1<u^2<u^3`$ where $`\nu `$ is the parameter entering in the CH equation (1). Then following Hayes and Whitham the equations can be written in Hamiltonian form with a local Poisson bracket of Dubrovin-Novikov type $$u_t^i=C^i(𝒖)u_x^i=A^{ij}\frac{h}{u^j}$$ where $$A^{ij}=g^{ii}\delta ^{ij}\frac{d}{dx}g^{ii}\mathrm{\Gamma }_{ik}^ju_x^k$$ (5) is the Hamiltonian operator and $`h`$ the Hamiltonian density. As pointed out by Dubrovin and Novikov, $`A^{ij}`$ defines a Hamiltonian operator if and only if $`g^{ii}=g^{ii}(𝒖)`$ is a flat non degenerate metric and $`\mathrm{\Gamma }_{ik}^j`$ are the Christoffel symbols of the corresponding Levi-Civita connection. We also find a second local compatible Hamiltonian structure which is obtained from the flat metric $`g^{ii}(𝒖)(u^i+\nu )`$ where $`\nu `$ is the constant in the CH equation. Therefore the nonlocal bi-Hamiltonian structure of the original CH equation averages to a local bi-Hamiltonian structure of Dubrovin Novikov type. A reciprocal transformation is a closed form which changes the independent variables of the equation and maps conservation laws into conservations laws, but it does not preserve the Poisson structure as shown by E.V. Ferapontov and M.V. Pavlov ,. The Camassa Holm equation can be transformed by a reciprocal transformation into the first negative flow of the KdV hierarchy (also known as AKNS equation ). An elegant treatment of the relations among positive and negative flows of the CH and KdV hierarchies can be found in . Let $`g_{ii}^{KdV}`$ and $`g_{ii}^{KdV}/\beta ^i`$ be the flat compatible metrics associated to the bi-Hamiltonian structure of the KdV modulation equations with respect to the usual Riemann invariants $`\beta ^1,\beta ^2,\beta ^3`$ as defined in . Then the reciprocal transformation is generated by the Casimir $`_0`$ of the Hamiltonian operator associated to the metric $`g_{ii}^{KdV}/\beta ^i`$. According to the results in , the reciprocal transformation maps the two KdV flat metrics to the CH metrics $$\frac{g_{ii}^{KdV}}{_0^2},\frac{g_{ii}^{KdV}}{_0^2\beta ^i}$$ which are not flat. The relation between the CH Riemann invariants and the KdV Riemann invariants is $`\beta ^i=1/(u^i+\nu )`$. The corresponding CH modulation equations are Hamiltonian with respect to two non-local operators of Mokhov-Ferapontov and Ferapontov type , which are of the form (5) plus a nonlocal tail. However from the Lagrangian averaging, we independently prove the existence of one local Hamiltonian structure. We show that the two metrics $$\frac{g_{ii}^{KdV}}{_0^2(\beta ^i)^2},\frac{g_{ii}^{KdV}}{_0^2(\beta ^i)^3}$$ are flat and define a flat pencil, that is, the CH modulation equations are bi-Hamiltonian with respect to two local Hamiltonian operators of the form (5). We remark that the two flat KdV metrics $`g_{ii}^{KdV}`$ and $`g_{ii}^{KdV}/\beta ^i`$ are related to a semisimple Frobenius manifold . More in general, B. Dubrovin proves that, under certain assumptions, a flat pencil of contravariant metrics on a manifold induces a Frobenius structure on it. One of the assumptions is the requirement that one of the two flat metrics is of Egorov type (namely its rotation coefficients are symmetric). Since none of the two CH flat metrics have the Egorov property, there is no Frobenius structure associated to this system. Therefore, from the geometric point of view, the KdV modulation equations and the CH modulation equation belong to two different classes. All the results presented here for the one-phase CH modulation equations may be generalized to the multi–phase case in a straightforward way. However, in the present paper we have decided to concentrate only on the one-phase case to better clarify similarities and differences with the KdV case, and we will present the discussion of the multi–phase case in a future publication. ## 2 Whitham modulation equations In this section we use Lagrangian formalism to average the Camassa-Holm equation in the genus one case and refer to , for a general exposition of the method we use. Introducing the potential $$\varphi :\varphi _x=𝚞,$$ equation (1) takes the form $$\varphi _{xt}\varphi _{xxxt}+3\varphi _x\varphi _{xx}2\varphi _{xx}\varphi _{xxx}\varphi _x\varphi _{xxxx}+2\nu \varphi _{xx}=0,$$ and a Lagrangian is $$=\frac{1}{2}\varphi _x\varphi _t+\frac{1}{2}\varphi _{xxx}\varphi _t\frac{1}{2}\varphi _x^3\nu \varphi _x^2+\frac{1}{4}\varphi _x^2\varphi _{xxx}.$$ (6) We consider $`2\pi `$–periodic solutions of the form $$𝚞=\eta (\theta ),\theta =kx\omega t.$$ Following Whitham , we introduce the pseudo-phase $$\varphi =\psi +\mathrm{\Phi }(\theta ),\psi =\beta x\gamma t,\theta =kx\omega t,$$ where $`\mathrm{\Phi }(\theta )`$ is a $`2\pi `$ periodic function of $`\theta `$ with zero average. The averaged Lagrangian over the one-dimensional real torus is $`\overline{}`$ $`={\displaystyle }d\theta [{\displaystyle \frac{1}{2}}(\beta +k\mathrm{\Phi }_\theta )(\gamma \omega \mathrm{\Phi }_\theta )+{\displaystyle \frac{1}{2}}k^3\mathrm{\Phi }_{\theta \theta \theta }(\gamma \omega \mathrm{\Phi }_\theta )`$ $`{\displaystyle \frac{1}{2}}(\beta +k\mathrm{\Phi }_\theta )^3\nu (\beta +k\mathrm{\Phi }_\theta )^2+{\displaystyle \frac{1}{4}}(\beta +k\mathrm{\Phi }_\theta )^2k^3\mathrm{\Phi }_{\theta \theta \theta }]`$ Now we suppose that the constants $`\beta `$, $`k`$, $`\gamma `$ and $`\omega `$ are slowly varying functions of time, that is $`\beta =\beta (X,T)`$, $`\gamma =\gamma (X,T)`$, $`k=k(X,T)`$ and $`\omega =\omega (X,T)`$ with $`X`$ and $`T`$ “slow variables”. Then $$\overline{}=\overline{}(k,\omega ,\beta ,\gamma ;X,T).$$ The equations of slow modulation of the parameters $`k,\omega ,\beta ,\gamma `$ are the extremals of the functional $$\overline{}(k,\omega ,\beta ,\gamma )𝑑X𝑑T,$$ and take the form $$\{\begin{array}{c}_X\overline{}_k+_T\overline{}_\omega =0,k_T+\omega _X=0,\\ _X\overline{}_\gamma +_T\overline{}_\beta =0,\beta _T+\gamma _X=0.\end{array}$$ (7) As in the KdV case, substituting the forth equation into the third in (7), the $`X`$-derivative disappears and we get a constraint. Thus the number of equations reduces from four to three. We get to the same conclusion transforming the modulation equations in Hamiltonian form. Following Hayes and Whitham , let us introduce the Hamiltonian density $$=(k,\omega ,\beta ,\gamma )=\omega \overline{}_\omega +\gamma \overline{}_\gamma \overline{}.$$ Then the modulation equations are Hamiltonian with respect to the canonical Poisson bracket $$\{k(X),\overline{}_\omega (Y)\}=\delta ^{}(XY),\{\beta (X),\overline{}_\gamma (Y)\}=\delta ^{}(XY).$$ Since the constraint $$\overline{}_\gamma =\frac{1}{2}\beta ,$$ the number of fields reduces from four to three. This is connected with the Dirac reduction. Therefore the Whitham equations can be written in Hamiltonian form with a local Dubrovin-Novikov Poisson bracket $$\{k(X),\overline{}_\omega (Y)\}=\delta ^{}(XY),\{\beta (X),\beta (Y)\}=2\delta ^{}(XY)$$ and with Hamiltonian $`=(k,J,\beta )`$, $`J=\overline{}_\omega `$. In these variables the equations of motion are Hamiltonian $$k_T=_X_J,J_T=_X_k,\beta _T=_X_\beta .$$ (8) ## 3 Modulation equations for CH in Riemann invariant form The modulation equations (8) can also be written in Riemann invariant form. For the purpose we introduce the spectral curve associated to the periodic travelling wave solution $`𝚞(x,t)=\eta (kx\omega t)`$. When we plug $`\eta (\theta )`$, $`\theta =kx\omega t`$, into the CH equation (1), we get, after integration, $$k^2(c\eta )\eta _\theta ^2+(2\nu c)\eta ^2+\eta ^3+2B\eta 2A=0,$$ (9) where $`A`$ and $`B`$ are constants of integration and $`c=\omega /k`$. The CH one-phase solution $`𝚞(x,t)=\eta (kx\omega t)`$ is obtained by inverting the third kind differential $$_{u_0}^u\frac{(\eta c)d\eta }{\sqrt{(\eta c)(\eta ^3+(2\nu c)\eta ^2+2B\eta 2A)}}=kx\omega t.$$ The inversion of the above integral is discussed in where the solution $`𝚞(x,t)`$ is expressed in terms of convenient generalized theta-functions in two variables, which are constrained to the generalized theta-divisor. Integration of (9) over $`\theta `$ yields the amplitude dependent dispersion relation for the nonlinear dispersive wave $$k\frac{(\eta c)d\eta }{\sqrt{(\eta c)(\eta ^3+(2\nu c)\eta ^2+2B\eta 2A)}}=2\pi ,$$ (10) where the integration is taken on a closed path between $`e^2`$ and $`e^1`$ where $`c>e^1>e^2>e^3`$ are the roots of the polynomial $$\eta ^3(c2\nu )\eta ^2+2B\eta 2A,$$ with the constraint $$2\nu =ce^1e^2e^3.$$ Here and below, we denote vectors with upper indices. ¿From now on, we will use small letters $`x,t`$ for the ‘slow variables’ $`X,T`$ introduced in the previous section, since we deal only with modulation equations and no ambiguity may occur. ###### Theorem 3.1 The one-phase CH modulation equations (8) take the Riemann invariant form $$_tu^i+C^i(𝒖)_xu^i=0,i=1,\mathrm{},3,$$ (11) where the Riemann invariants $`𝐮=(u^1,u^2,u^3)`$, $`u^1<u^2<u^3`$, are $$u^1=\frac{1}{2}(e^2+e^3),u^2=\frac{1}{2}(e^1+e^3),u^3=\frac{1}{2}(e^1+e^2),$$ and the speeds $`C^i(𝐮)`$ take the form $$C^i(𝒖)=\frac{_{u^i}\omega (𝒖)}{_{u^i}k(𝒖)}.$$ (12) The wave number and frequency are given by the Abelian integrals $$\begin{array}{cc}\hfill k& =2\pi \left(_a\frac{(\lambda +\nu )d\lambda }{\sqrt{R(\lambda )}}\right)^1,R(\lambda )=(\lambda +\nu )(\lambda u^1)(\lambda u^2)(\lambda u^3),\hfill \\ \hfill \omega & =(2\nu +u^1+u^2+u^3)k,\hfill \end{array}$$ (13) and the integration is taken on cycle $`a`$ passing between $`u^2`$ and $`u^1`$ (see figure 1). The existence of the Riemann invariants and equations (12) is proven directly starting from (8) and using variational identities of Abelian integrals. To write the velocities in an explicit form, we introduce the integrals $$I_k=_a\frac{\lambda ^k}{\sqrt{R(\lambda )}}𝑑\lambda ,k0,$$ (14) where $`R(\lambda )`$ is defined in (13). Let $`\sigma _1(\lambda )`$ be the normalized third kind differential with first order pole at $`(\mathrm{},\pm \mathrm{})`$ with residue $`\pm 1`$, respectively, and let $`\sigma _2`$ be the normalized second kind differential with second order pole at infinity, namely $`\sigma _1(\lambda )={\displaystyle \frac{P_1(\lambda )d\lambda }{\sqrt{R(\lambda )}}},P_1(\lambda )=\lambda +\gamma _1,`$ (15) $`\sigma _2(\lambda )={\displaystyle \frac{P_2(\lambda )d\lambda }{\sqrt{R(\lambda )}}},P_2(\lambda )=\lambda ^2{\displaystyle \frac{1}{2}}(u^1+u^2+u^3\nu )\lambda +\gamma _2`$ (16) where the constants $`\gamma _1={\displaystyle \frac{I_1}{I_0}}`$ and $`\gamma _2={\displaystyle \frac{I_2}{I_0}}+{\displaystyle \frac{1}{2}}(u^1+u^2+u^3\nu ){\displaystyle \frac{I_1}{I_0}}`$ are uniquely determined by $$_a\sigma _i(\lambda )=0,i=1,2.$$ (17) These constants are explicitly given by $$\gamma _2=\frac{1}{2}[u^1u^2\nu u^3+(u^3u^1)(u^2+\nu )\frac{E(s)}{K(s)}],\gamma _1=\nu (u^1+\nu )\frac{\mathrm{\Lambda }(K(s),\rho ,s)}{K(s)},$$ (18) where $$K(s)=_0^{\pi /2}\frac{d\psi }{\sqrt{1s^2\mathrm{sin}^2\psi }},E(s)=_0^{\pi /2}𝑑\psi \sqrt{1s^2\mathrm{sin}^2\psi }$$ are the complete elliptic integrals of the first and second kind respectively with modulus $$s^2=\frac{(u^2u^1)(u^3+\nu )}{(u^3u^1)(u^2+\nu )}$$ and $$\mathrm{\Lambda }(K(s),\rho ,s)=_0^{K(s)}\frac{dv}{1\rho ^2sn^2v},\rho ^2=\frac{u^2u^1}{u^2+\nu },$$ is the complete elliptic integral of the third kind with $`sn`$ the Jacobi elliptic function. ###### Theorem 3.2 The speeds $`C^i(𝐮)`$, $`i=1,2,3`$ defined in (12) take the form $$\begin{array}{cc}& C^1(u^1,u^2,u^3)=u^1+u^2+u^3+2\nu +2\frac{(u^1+\nu )(u^1u^2)\mathrm{\Lambda }(K(s),\rho ,s)}{(u^2+\nu )[K(s)E(s)]}\hfill \\ & C^2(u^1,u^2,u^3)=u^1+u^2+u^3+2\nu +\frac{2(u^2u^1)\mathrm{\Lambda }(K(s),\rho ,s)}{K(s){\displaystyle \frac{(u^2+\nu )(u^3u^1)}{(u^1+\nu )(u^3u^2)}}E(s)}\hfill \\ & C^3(u^1,u^2,u^3)=u^1+u^2+u^3+2\nu +2\frac{(u^1+\nu )(u^3u^2)\mathrm{\Lambda }(K(s),\rho ,s)}{(u^2+\nu )E(s)}.\hfill \end{array}$$ (19) where $`K(s)`$, $`E(s)`$ and $`\mathrm{\Lambda }(K(s),\rho ,s)`$ are the complete elliptic integrals of first, second and third kind with modulus $`s^2={\displaystyle \frac{(u^2u^1)(u^3+\nu )}{(u^3u^1)(u^2+\nu )}}`$. The equations $`_tu^i+C^i(𝒖)_xu^i=0`$ are hyperbolic and the velocities satisfy $$C^1(𝒖)<C^3(𝒖),C^2(𝒖)<C^3(𝒖),\nu <u^1<u^2<u^3.$$ In the limit when two Riemann invariants coalesce, the modulation equation reduce to the dispersionless CH equation $$_t𝚞+(3𝚞+2\nu )_x𝚞=0.$$ To prove the theorem we introduce the normalized holomorphic differential $`\varphi (\lambda )`$ $$\varphi (\lambda )=\frac{d\lambda }{I_0\sqrt{R(\lambda )}},_a\varphi (\lambda )=1.$$ Next we observe that the wave number $`k`$ defined in (13) takes the form $$k=2\pi \frac{\varphi (\nu )}{\sigma _1(\nu )},$$ (20) where $$\sigma _1(\nu ):=\frac{2P_1(\nu )}{\sqrt{(\nu u^1)(\nu u^2)(\nu u^3)}}=\frac{\sigma _1(\lambda )}{dt}\|_{\lambda =\nu },t^2=\lambda +\nu ,$$ and $$\varphi (\nu ):=\frac{2}{I_0\sqrt{(\nu u^1)(\nu u^2)(\nu u^3)}}=\frac{\varphi (\lambda )}{dt}\|_{\lambda =\nu },t^2=\lambda +\nu ,$$ with $`\sigma _1(\lambda )`$ defined in (15). The following variational formulas hold $$\frac{}{u^i}\varphi (\nu )=\frac{1}{2}\varphi (u^i)\mathrm{\Omega }_\nu (u^i),\frac{}{u^i}\sigma _1(\nu )=\frac{1}{2}\sigma _1(u^i)\mathrm{\Omega }_\nu (u^i),$$ (21) where $$\varphi (u^i):=\frac{\varphi (\lambda )}{dt}\|_{\lambda =u^i},\mathrm{\Omega }_\nu (u^i):=\frac{\mathrm{\Omega }_\nu (\lambda )}{dt}\|_{\lambda =u^i},t^2=\lambda u^i.$$ (22) and $`\mathrm{\Omega }_\nu (\lambda )`$ is a second kind Abelian differential with second order pole at $`\lambda =\nu `$ with asymptotic behavior for $`\lambda \nu `$ $$\mathrm{\Omega }_\nu (\lambda )\frac{dt}{t^2},t^2=\lambda +\nu ,$$ and normalized by the condition $$_a\mathrm{\Omega }_\nu (\lambda )=0.$$ (23) The differential $`\mathrm{\Omega }_\nu (\lambda )`$ is explicitly given by the expression $$\begin{array}{cc}\hfill \mathrm{\Omega }_\nu (\lambda )& =\frac{P_\nu (\lambda )d\lambda }{\sqrt{R(\lambda )}\sqrt{(\nu u^1)(\nu u^2)(\nu u^3)}},\hfill \\ \hfill P_\nu (\lambda )& =\frac{(\nu +u^1)(\nu +u^2)(\nu +u^3)}{2(\lambda +\nu )}+P_2(\nu ),\hfill \end{array}$$ (24) where $`P_2(\lambda )`$ has been defined in (16). Inserting (20) and (21) into (12), we finally obtain the expression $$C^i(𝒖)=u^1+u^2+u^3+2\nu \frac{P_1(\nu )}{P_\nu (u^i)}\underset{ji,j=1}{\overset{3}{}}(u^iu^j),i=1,2,3,$$ (25) where the rational function $`P_1(\lambda )`$ and $`P_\nu (\lambda )`$ are as in (15) and (24) respectively. Finally inserting (18) into (25) we arrive at the formula (19). Next we prove the ordering of the velocities using the formula (25). Since $`\nu <u^1<u^2<u^3`$, and $`\underset{\lambda \nu ^+}{lim}P_\nu (\lambda )=\mathrm{}`$, by monotonicity, there is only one point $`\lambda ^{}>\nu `$ for which $$\frac{(\nu +u^1)(\nu +u^2)(\nu +u^3)}{2(\lambda ^{}+\nu )}=P_2(\nu );$$ moreover, because of (23), the point $`\lambda ^{}`$ satisfies the inequality $`u^1<\lambda ^{}<u^2`$ . Therefore $$P_\nu (u^1)<0,P_\nu (u^3)>P_\nu (u^2)>0.$$ In the same way, using (17) we conclude that $`P_1(\nu )<0`$. Using the above inequalities, it is straightforward to verify that $$C^1(𝒖)<C^3(𝒖),,C^2(𝒖)<C^3(𝒖),\nu <u^1<u^2<u^3.$$ In general there is not a strict ordering between $`C^1(𝒖)`$ and $`C^2(𝒖)`$. For example for step-like initial data the characteristics $`C^1(𝒖)`$ and $`C^2(𝒖)`$ do cross for $`u^1<u^2<u^3`$. We end the proof of the theorem studying the behavior of the speeds $`C^i(𝒖)`$, $`i=1,2,3`$, when two of the Riemann invariants coalesce. In the limiting case $$u^2=vϵ,u^3=v+ϵ,ϵ0,$$ we have $$\underset{ϵ0}{lim}P_1(\lambda )=\lambda v,\underset{ϵ0}{lim}P_\nu (\lambda )=\frac{1}{2(\lambda +\nu )}(\lambda v)(u^1+\nu )(v+\nu ),$$ so that $$\underset{ϵ0}{lim}\frac{P_1(\nu )}{P_\nu (u^1)}(u^1u^2)(u^1u^3)=2(vu^1)$$ and $$\underset{ϵ0}{lim}\frac{P_1(\nu )}{P_\nu (u^2)}(u^2u^1)(u^2u^3)=\underset{ϵ0}{lim}\frac{P_1(\nu )}{P_\nu (u^3)}(u^3u^2)(u^3u^1).$$ Combining the above relations, we have $$C^2(u^1,v,v)=C^3(u^1,v,v),C^1(u^1,v,v)=3u^1+2\nu ,$$ which gives the dispersionless CH equation $`_tu^1+(3u^1+2\nu )_xu^1=0`$. In the same way, it can be proved that, in the limit $`u^2=u^1`$, the speeds $`C^1(u^1,u^1,u^3)=C^2(u^1,u^1,u^3)`$ and $`C^3(u^1,u^1,u^3)=3u^3+2\nu `$. ### 3.1 Hamiltonian structure and integration In this section, we investigate the bi-Hamiltonian structure of the one-phase Whitham equations $$_tu^i+C^i(𝒖)_xu^i=0,i=1,\mathrm{},3,$$ with $`C^i(𝒖)`$ as in (12) or (25). In section 2 we have proven that the above equations are Hamiltonian with respect to a canonical Poisson bracket. In the following we show that the CH modulation equations are bi-Hamiltonian with respect to local Poisson brackets of Dubrovin-Novikov type. ###### Proposition 3.3 The speeds $`C^i(𝐮)`$, $`i=1,2,3,`$ satisfy the following relations $$\frac{_{u^j}C^i}{C^jC^i}=_{u^j}\mathrm{log}\sqrt{g_{ii}},$$ (26) where $`g_{ii}=g_{ii}(𝐮)`$ is the covariant metric defined by the relation $$g_{ii}(𝒖)=4(u^i+\nu )\frac{\underset{[}{Res}\lambda =u^i]\{{\displaystyle \frac{(\mathrm{\Omega }_\nu (\lambda ))^2}{d\lambda }}\}}{\underset{[}{Res}\lambda =\nu ]\{{\displaystyle \frac{(\sigma _1(\lambda ))^2}{d\lambda }}\}},$$ (27) with $`\mathrm{\Omega }_\nu (\lambda )`$ the second kind differential defined in (24) and $`\sigma _1(\lambda )`$ the third kind differential defined in (15). The metric $`g_{ii}(𝐮)`$ is flat. The metric $`g_{ii}(𝐮)`$ is defined up to multiplication by an arbitrary function $`f_i(u^i)`$. The metrics $$\frac{g_{ii}(𝒖)}{2(u^i+\nu )},\frac{g_{ii}(𝒖)}{4(u^i+\nu )^2},\frac{g_{ii}(𝒖)}{8(u^i+\nu )^3},$$ (28) are respectively flat, of constant curvature $`R_{ij}^{ij}=1`$ and conformally flat with curvature tensor $`R_{ij}^{ij}=C^i(𝐮)C^j(𝐮)2\nu `$. The pencil of metrics $$g_{ii}(𝒖)+\lambda \frac{g_{ii}(𝒖)}{(u^i+\nu )}$$ (29) is flat for any real $`\lambda `$. None of the metrics defined in (27) and in (28) is of Egorov type. We denote the diagonal metric in covariant form by $`g_{ii}`$ and its inverse by $`g^{ii}`$. To derive (26) it is sufficient to use the variational formulas (21) and the additional one $$\frac{}{u^i}\mathrm{\Omega }_\nu (u^j)=\frac{1}{2}\mathrm{\Omega }_\nu (u^i)\mathrm{\Omega }_{u^i}(u^j),$$ where $`\mathrm{\Omega }_{u^i}(\lambda )`$ is a normalized second kind differential with second order pole at $`\lambda =u^i`$. The explicit form of $`\mathrm{\Omega }_{u^i}(\lambda )`$ can be obtained from (24) by replacing $`\nu `$ by $`u^i`$. The quantities $`\mathrm{\Omega }_\nu (u^i)`$ and $`\mathrm{\Omega }_{u^i}(u^j)`$ are defined as in (22). To prove the second part of the proposition we evaluate the non-zero elements of the curvature tensor $`R_{il}^{ij}`$ $`R_{il}^{ij}`$ $`={\displaystyle \frac{1}{\sqrt{g_{ii}g_{jj}}}}\left\{_lr_{ji}r_{jl}r_{li}\right\},ijl,`$ (30) $`R_{ij}^{ij}`$ $`={\displaystyle \frac{1}{\sqrt{g_{ii}g_{jj}}}}\left\{_ir_{ij}+_jr_{ji}+{\displaystyle \underset{pi,j}{}}r_{pj}r_{pi}\right\}.`$ (31) Here $`r_{ij}`$ are the rotation coefficients defined by $$r_{ij}=\frac{_{u^i}\sqrt{g_{jj}}}{\sqrt{g_{ii}}},ij.$$ (32) A metric is of Egorov type if $`r_{ij}=r_{ji}`$, $`ij`$. By direct calculation we obtain $$r_{ij}=\frac{1}{2}\left(\frac{u^j+\nu }{u^i+\nu }\right)^{k/2}\left(\mathrm{\Omega }_{u^i}(u^j)\frac{\mathrm{\Omega }_\nu (u^j)\sigma _1(u^i)}{\sigma _1(\nu )}\right),k=1,0,1,2.$$ ¿From the above formula it is clear that $`r_{ij}r_{ji}`$, for all the four metrics and therefore none of them is an Egorov metric. To evaluate (30-31) the following additional variational formulas are needed $$\begin{array}{cc}\hfill \frac{}{u^i}\gamma _1& =\frac{1}{2}+\frac{1}{4}\sigma _1(u^i)\sigma _2(u^i),i=1,2,3,\hfill \\ \hfill \frac{}{u^i}\gamma _2& =\frac{1}{4}(u^1+u^2+u^3\nu )\frac{1}{2}u^i+\frac{1}{4}(\sigma _2(u^i))^2,i=1,2,3,\hfill \end{array}$$ (33) where $`\sigma _1`$ and $`\sigma _2`$ have been defined in (15) and (16) respectively and $$\sigma _k(u^i):=\frac{\sigma _k(\lambda )}{dt}\|_{\lambda =u^i},t^2=\lambda u^i,i=1,2,3,k=1,2.$$ Using the variational formulas (33) we obtain that $`R_{il}^{ij}=0`$ for all the four metrics and $$\begin{array}{cc}\hfill g_{ii}(𝒖)& R_{ij}^{ij}=0,\hfill \\ \hfill \frac{g_{ii}(𝒖)}{2(u^i+\nu )}& R_{ij}^{ij}=0,\hfill \\ \hfill \frac{g_{ii}(𝒖)}{4(u^i+\nu )^2}& R_{ij}^{ij}=1,\hfill \\ \hfill \frac{g_{ii}(𝒖)}{8(u^i+\nu )^3}& R_{ij}^{ij}=2\nu C^i(𝒖)C^j(𝒖).\hfill \end{array}$$ The flatness of the pencil of metrics (29) can be obtained from the results in . An elegant proof of the flatness of the metrics, valid for any genus can be obtained in a more convenient set of coordinates $`u^i{\displaystyle \frac{1}{u^i+\nu }}`$ that makes the spectral curve odd with a branch point at infinity. We use these coordinates in the next section to compare the Whitham equations for CH and for KdV. ###### Remark 3.4 The non-existence of a flat Egorov metric is related to the non-existence of conservation laws of the form $`a_t=b_x`$ and $`b_t=c_x`$ . Furthermore the non-existence of a flat Egorov metric implies that the CH modulation equations cannot be associated to a Frobenius manifold. We recall that, under certain assumptions, a flat pencil of contravariant metrics on a manifold induces a Frobenius structure on it . One of the assumptions is the requirement that one of the two flat metrics is an Egorov metric. Therefore the geometric structure of the CH modulation equations is substantially different from that of the KdV modulation equations. To any flat diagonal Riemannian controvariant metric $`g^{ii}`$, Dubrovin and Novikov associate a local homogenous Poisson bracket of hydrodynamic type $$\{F,G\}=\frac{\delta F}{\delta u^i}A^{ij}\frac{\delta G}{\delta u^j}$$ defined by the Hamiltonian operator $$A^{ij}=g^{ii}\delta ^{ij}\frac{d}{dx}g^{ii}\mathrm{\Gamma }_{ik}^ju_x^k.$$ (34) Indeed in , they prove that $`A^{ij}`$ defines a Poisson tensor if and only if $`g^{ii}`$ is a flat non-degenerate metric and $`\mathrm{\Gamma }_{ik}^j`$ are the Christoffel symbols of the Levi-Civita connection compatible with the metric (the metric is not necessary diagonal in their formulation). If the metric $`g^{ii}`$ is not flat, the Poisson tensor needs to be modified adding a non-local tail , of the form $$A^{ij}=g^{ii}\delta ^{ij}\frac{d}{dx}g^{ii}\mathrm{\Gamma }_{ik}^ju_x^k+cu_x^i\left(\frac{d}{dx}\right)^1u_x^j,$$ (35) for metrics of constant curvature $`c`$, and of the form $$A^{ij}=g^{ii}\delta ^{ij}\frac{d}{dx}g^{ii}\mathrm{\Gamma }_{ik}^ju_x^k+\eta ^iu_x^i\left(\frac{d}{dx}\right)^1u_x^j+u_x^i\left(\frac{d}{dx}\right)^1\eta ^ju_x^j,$$ (36) for conformally flat metrics. In the above relation the affinors $`\eta ^j`$ satisfy the equations $$R_{ij}^{ij}=\eta ^i+\eta ^j,$$ where $`R_{ij}^{ij}`$ is the curvature tensor. In the following let $`g_{ii}(𝒖)`$ be as in (27). Let $`A_1^{ij}`$ and $`A_2^{ij}`$ be the Hamiltonian operators of the form (34) that correspond to the flat metric $`g_{ii}(𝒖)`$ and to $`{\displaystyle \frac{g_{ii}(𝒖)}{2(u^i+\nu )}}`$, respectively. The linear combination $`A_1^{ij}+\lambda A_2^{ij}`$ is an Hamiltonian operator for any $`\lambda `$ because $`g_{ii}(𝒖)+\lambda g_{ii}(𝒖)/u^i`$ is a flat pencil of metrics for any $`\lambda `$. The Hamiltonian structure obtained in (8) coincides with Hamiltonian operator $`A_1`$. Indeed the coordinates $`k=k(u^1,u^2,u^3)`$, $`\beta =\beta (u^1,u^2,u^3)`$ and $`J=J(u^1,u^2,u^3)`$ defined in (8) are the densities of the Casimirs for the Hamiltonian operator $`A_1^{ij},`$ namely $$A_1^{ij}\frac{\delta k}{\delta u^j}=A_1^{ij}\frac{\delta \beta }{\delta u^j}=A_1^{ij}\frac{\delta J}{\delta u^j}=0,$$ and therefore they are the flat coordinates for the metric $`g_{ii}(𝒖)`$ . Defining $$h_0:=m𝑑\theta =\beta =2\frac{P_2(\nu )}{P_1(\nu )}\nu ,$$ where $`P_1`$ and $`P_2`$ have been defined in (15) and (16) respectively, we obtain the Hamiltonian densities $`h_k=h_k(𝒖)`$, $`k0`$, by the recursion relation $$A_2^{ij}\frac{h_k}{u^j}=A_1^{ij}\frac{h_{k+1}}{u^j},k0.$$ The CH modulation hierarchy takes the bi-Hamiltonian form $$_{t_k}u^i=A_1^{ij}\frac{\delta h_{k+1}}{\delta u^j}=A_2^{ij}\frac{\delta h_k}{\delta u^j},i=1,2,3.$$ ###### Theorem 3.5 For $`k=1`$ and $`t_1=t`$, the above equations coincide with the modulation equations obtained in (11) and take the bi-Hamiltonian form $$u_t^i=A_1^{ij}\frac{\delta h_2}{\delta u^j}=A_2^{ij}\frac{\delta h_1}{\delta u^j}$$ where the Hamiltonian densities are the averaged Hamiltonians $$h_2=\frac{1}{2}(𝚞^3+\mathrm{𝚞𝚞}_x^2+2\nu 𝚞^2)𝑑\theta ,h_1=\frac{1}{2}(𝚞^2+𝚞_x^2)𝑑\theta .$$ Moreover, the generating function for the Hamiltonian densities $`h_k`$, $`k0`$, is given by the coefficients of the expansion as $`\lambda \mathrm{}`$ of the differential $$\frac{\mathrm{\Omega }_\nu (\lambda )}{\{\underset{[}{Res}\lambda =\nu ]{\displaystyle \frac{(\sigma _1(\lambda ))^2}{d\lambda }}\}^{\frac{1}{2}}}(\xi _0+\xi _1\frac{1}{\lambda }+\xi _2\frac{1}{\lambda ^2}+\mathrm{})\frac{d\lambda }{\lambda ^2}.$$ In particular, the first Hamiltonians (modulo Casimirs) are $$h_0=2\xi _0\nu =\beta ,h_1=2\xi _1+2\nu \xi _0,h_2=\frac{8}{3}\xi _2+6\nu \xi _1.$$ (37) In the following we write the CH-modulation equation in Hamiltonian form with respect to a nonlocal Hamiltonian operator of Mokhov-Ferapontov and Ferapontov type. Let $`A_3^{ij}`$ be the Hamiltonian operator of the form (35) associated to the metric $`{\displaystyle \frac{g_{ii}(𝒖)}{4(\nu +u^i)^2}}`$ of constant curvature $`c=1`$, and let $`A_4^{ij}`$ be the Hamiltonian operator of the form (36) associated to the conformally flat metric $`{\displaystyle \frac{g_{ii}(𝒖)}{8(\nu +u^i)^3}}`$ with affinors $`\eta ^i=C^i+\nu `$, $`i=1,2,3`$. The CH modulation equations (11) can also be written in a non-local Hamiltonian form $$_tu^i=A_3^{ij}\frac{\delta h_0}{\delta u^j}=A_4^{ij}\frac{\delta h_1}{\delta u^j},i=1,2,3,$$ where $$h_1=1\frac{\nu }{\{\underset{[}{Res}\lambda =\nu ]{\displaystyle \frac{(\sigma _1(\lambda ))^2}{d\lambda }}\}^{\frac{1}{2}}}.$$ We remark that the Hamiltonian operators $`A_3`$ and $`A_4`$ can be obtained from the recursion $`A_3=^2A_1`$ and $`A_4=^3A_1`$ where $`=A_2A_1^1`$ is the recursion operator and $`A_1^1`$ denotes the inverse of $`A_1`$. In the limit when two Riemann invariants coalesce, the Hamiltonian operators $`A_1^{ij}`$ and $`A_2^{ij}`$ reduce to the “dispersionless limit” of the Poisson operators of the CH equation $$𝒫_1=\frac{d}{dx},𝒫_2=2(𝚞+\nu )\frac{d}{dx}𝚞_x,$$ respectively. In the same limit, the Hamiltonian operators $`A_3^{ij}`$ and $`A_4^{ij}`$ reduce to $`𝒫_2𝒫_1^1𝒫_2`$ and $`𝒫_2𝒫_1^1𝒫_2𝒫_1^1𝒫_2`$, respectively. ### 3.2 Integration of the Whitham equations In this subsection we show how to integrate the Whitham equations (11). All hydrodynamic systems satisfying (26) are integrable, via the generalized hodograph transform introduced by Tsarev . Indeed such systems (not necessarily with local Hamiltonian) possess an infinite number of commuting flows $$_t^{}u^i=w^i(𝒖)_xu^i,$$ where the $`w^i`$ are solutions of the linear overdetermined system $$\frac{_{u^j}w^i}{w^iw^j}=\frac{_{u^j}C^i}{C^iC^j},ij,$$ (38) where $`C^i=C^i(u^1,u^2,u^3)`$, $`i=1,2,3,`$ are the speeds in (12). Then the solution $`𝒖(x,t)=(u^1(x,t),u^2(x,t),u^3(x,t))`$ of the so-called hodograph transform $$x=C^i(𝒖)t+w^i(𝒖)i=1,2,3,$$ (39) satisfies the system (11). Conversely, any solution $`(u^1(x,t),u^2(x,t),u^3(x,t))`$ of (11) can be obtained in this way. For monotone decreasing initial data $`x=f(u)\|_{t=0}`$, the general solution of the system (38) can be obtained following the work of Fei-Ran Tian and the algebraic-geometric integration of Krichever . For simplicity, we restrict ourselves to the case $`\nu =0`$. ###### Proposition 3.6 For $`\nu =0`$, and monotone increasing initial data $`x=f(u)\|_{t=0}`$, the solution of the system (38) is $$w^i(𝒖)=q(𝒖)+\left(C^i(𝒖)u^1u^2u^3\right)\frac{q(𝒖)}{u^i},$$ (40) where the function $`q=q(𝐮)`$ solves the linear over-determined system of Euler-Poisson-Darboux type $$\begin{array}{cc}& _{u^i}q(𝒖)_{u^j}q(𝒖)=2(u^iu^j)_{u^i}_{u^j}q(𝒖),\hfill \\ & q(u,u,u)=f(u).\hfill \end{array}$$ (41) The proof of the above statement follows from . Equation (LABEL:euler) can be integrated and the explicit expression of the function $`q(u^1,u^2,u^3)`$ is $$q(u^1,u^2,u^3)=\frac{1}{2\sqrt{2\pi }}_1^1_1^1\frac{f\left({\displaystyle \frac{1+\mu }{2}}{\displaystyle \frac{1+\eta }{2}}u^1+{\displaystyle \frac{1+\mu }{2}}{\displaystyle \frac{1\eta }{2}}u^2+{\displaystyle \frac{1\mu }{2}}u^3\right)}{\sqrt{(1\mu )(1\eta ^2)}}𝑑\mu 𝑑\eta .$$ ## 4 CH modulation equations versus KdV and reciprocal transformations In this section we compare CH and KdV modulation equations. We start recalling the reciprocal transformation which links the CH equation to the first negative KdV flow. We show that the CH modulation equations (11) are transformed to the modulation equations of the first negative KdV flow by the averaged reciprocal transformation. Finally we compare the averaged Hamiltonian operators of the two systems. ### 4.1 A reciprocal transformation between Camassa Holm equation and the first negative flow of KdV In this subsection we summarize the relation between the Camassa Holm equation and the first negative flow of KdV hierarchy . The (associated) Camassa-Holm equation is transformed into the first negative KdV flow by a reciprocal transformation. In the following, to distinguish between CH and KdV, we use $`(x,t)`$ for Camassa-Holm variables and $`(y,\tau _{})`$ for the KdV variables. The change of dependent variable $`\rho ^2=m+\nu `$ transforms the Camassa–Holm equation $$m_t=2m𝚞_x𝚞m_x2\nu 𝚞_x,m=𝚞𝚞_{xx}.$$ into the associated Camassa-Holm equation $$\{\begin{array}{c}\rho _t=\left(𝚞\rho \right)_x,\hfill \\ \rho ^2=𝚞𝚞_{xx}+\nu ,\hfill \end{array}$$ (42) which, via the reciprocal transformation introduced by Fuchssteiner , $$\{\begin{array}{c}dy=\rho dx𝚞\rho dt,\hfill \\ d\tau _{}=dt,\hfill \end{array}$$ (43) is finally transformed into $$\{\begin{array}{c}𝚞=\rho ^2\nu \rho _{y\tau _{}}+\frac{\rho _\tau _{}\rho _y}{\rho },\hfill \\ \left(\frac{1}{\rho }\right)_\tau _{}=2\rho \rho _y\left(\rho \left(\mathrm{log}\rho \right)_{y\tau _{}}\right)_y.\hfill \end{array}$$ (44) The transformation (43) is a reciprocal transformation because the one-form $`\rho dx𝚞\rho dt`$ is closed with respect to the CH flow. The equation (44) is equivalent to the first negative flow of the KdV hierarchy $$\left(_y^2+2U+U_y_y^1\right)U_\tau _{}=0,$$ (45) under the condition $`U_\tau _{}=2\rho _y`$. Equation (45) may be re-expressed as $$\{\begin{array}{c}U=\frac{\rho _y^22\rho \rho _{yy}1}{2\rho ^2},\hfill \\ \left(\frac{\rho _y^22\rho \rho _{yy}1}{4\rho ^2}\right)_\tau _{}=\rho _y.\hfill \end{array}$$ (46) Finally, we observe that $`{\displaystyle \rho (x,t)𝑑x}`$ is a Casimir of the second Hamiltonian operator of the Camassa–Holm equation, $`P_2=m_x+_xm+2\nu _x`$. The 2$`\pi `$-periodic solutions of the first negative KdV flow (45), $`U(\mathrm{\Theta })`$, $`\mathrm{\Theta }=𝒦y\mathrm{\Omega }\tau `$, satisfy $$d\mathrm{\Theta }=\frac{𝒦dU}{\sqrt{U^3+\alpha U^2\beta U+\gamma }}.$$ (47) We may express such solutions also in the form $`\rho (\mathrm{\Theta })`$ or $`𝚞(\mathrm{\Theta })`$ and we easily get $$\{\begin{array}{c}𝚞_\mathrm{\Theta }^2=\frac{1}{C^2𝒦^2}(𝚞c)(𝚞^3(c2\nu )𝚞^2+2B𝚞2A).\hfill \\ \rho (\mathrm{\Theta })=\frac{C}{𝚞(\mathrm{\Theta })c},\hfill \\ U(\mathrm{\Theta })=\frac{2}{𝚞(\mathrm{\Theta })c},\hfill \end{array}$$ (48) where $`C=\mathrm{\Omega }/𝒦`$ and $`A,B,c`$ are the constants defined in (9) which satisfy $$Bc+\nu c^2A=C^2,\alpha =\frac{c^2+2B+4\nu c}{C^2}.$$ (49) The periodic solutions $`u(\theta )`$, $`\rho (\theta )`$, $`\theta =kx\omega t`$ of the (associated) Camassa–Holm equation, satisfy $$d\theta =\frac{k(𝚞c)d𝚞}{\sqrt{(𝚞c)(𝚞^3(c2\nu )𝚞^2+2B𝚞2A})}.$$ We observe that the reciprocal transformation sends 2$`\pi `$-periodic solutions in $`\mathrm{\Theta }`$ into 2$`\pi `$-periodic solution in $`\theta `$ (and vice versa). Indeed, let $`T`$ be the period of $`𝚞(\theta )`$, then $$T=\frac{k}{𝒦}_0^{2\pi }\frac{c𝚞}{C}𝑑\mathrm{\Theta }=k\frac{(𝚞c)d𝚞}{\sqrt{(𝚞c)(𝚞^3(c2\nu )𝚞^2+2B𝚞2A})}=2\pi .$$ It is then natural to expect that the average of the reciprocal transformation connects the CH modulation equations to the modulation equations of the first negative KdV flow. ### 4.2 The modulation equations of the negative KdV flow In this subsection, we compute the modulation equations of the KdV negative flow in the Riemann invariant coordinates used in the literature, namely the branch points $`\beta ^1,\beta ^2,\beta ^3`$, of the odd elliptic curve $$w^2=(\eta \beta ^1)(\eta \beta ^2)(\eta \beta ^3),$$ (50) where $$\beta ^1+\beta ^2+\beta ^3=\alpha ,$$ with $`\alpha `$ defined in (47). It turns out from (49) that $$\beta ^i=\frac{1}{u^i+\nu }>0,i=1,2,3,$$ (51) where $`u^i`$ are the CH Riemann invariants defined in theorem 3.1. On the Riemann surface (50) we define the second kind normalized differential $$dp(\eta )=\frac{\eta +\alpha _1}{\sqrt{(\eta \beta ^1)(\eta \beta ^2)(\eta \beta ^3)}}d\eta ,$$ (52) where $`\alpha _1`$ is uniquely determined by the normalization condition $$_a𝑑p(\eta )=0.$$ In the KdV literature $`dp(\eta )`$ is known as quasi-momentum. Now we can compute the velocities of the modulation equations of the negative KdV flow ###### Proposition 4.1 The one phase Whitham equations of the first negative KdV flow are $$_\tau _{}\beta ^i+v^i(\beta )_y\beta ^i=0,$$ (53) where $$v^i(\beta ):=\frac{_i\mathrm{\Omega }(\beta )}{_i𝒦(\beta )}=\frac{2}{\sqrt{\beta ^1\beta ^2\beta ^3}}\left(1\frac{_{ji}(\beta ^i\beta ^j)}{\beta ^i(\beta ^i+\alpha _1)}\right),$$ (54) with $`\alpha _1`$ as in (52). In the above relations $`\mathrm{\Omega }`$ and $`𝒦`$ are the frequency and wave-number of the one-phase KdV negative flow. The proof of the proposition is as follows. From (47) and (48), we immediately get $`𝒦=2\pi 𝒥_0^1`$ and $`\mathrm{\Omega }=C𝒦=4\pi 𝒥_0^1(\beta ^1\beta ^2\beta ^3)^{1/2}`$, where $$𝒥_0=\frac{d\lambda }{\sqrt{(\lambda \beta ^1)(\lambda \beta ^2)(\lambda \beta ^3)}}.$$ Then the expressions for the velocities (54) are computed from the definition using the following variational formula $$\frac{𝒥_0}{\beta ^i}=\frac{1}{2}𝒥_0\frac{\beta ^i+\alpha _1}{_{ji}(\beta ^i\beta ^j)},i=1,2,3.$$ (55) ###### Remark 4.2 In the limit $`\beta ^2=\beta ^3`$, the equations (53) converge to $`\beta _\tau _{}^1=2(\beta ^1)^{\frac{3}{2}}\beta _y^1`$, which is the dispersionless limit of the first negative KdV flow. Next, we write the negative KdV equations in Hamiltonian form. In the $`\beta `$s’ coordinates, the flat metric associated to the first local KdV-Whitham Hamiltonian operator is $$g_{ii}^{KdV}(𝜷)=\underset{[}{Res}\eta =\beta ^i]\{\frac{(dp(\eta ))^2}{d\eta }\},i=1,2,3,$$ (56) where the differential $`dp`$ has been defined in (52). ###### Remark 4.3 The flat metrics $`g_{ii}^{KdV}(𝛃)`$ and $`g_{ii}^{KdV}(𝛃)/\beta ^i`$ can be related to a Frobenius manifold defined on the moduli space of elliptic curves $`w^2=(\eta \beta ^1)(\eta \beta ^2)(\eta \beta ^3)`$ . Let $`J_1`$ and $`J_2`$ be the local KdV Hamiltonian operators of the form (34) associated to the flat metrics $`{\displaystyle \frac{1}{8}}g_{ii}^{KdV}(𝜷)`$ and $`{\displaystyle \frac{g_{ii}^{KdV}(𝜷)}{4\beta ^i}}`$. Let $`J_3`$ be the nonlocal Hamiltonian operator of the form (35) associated to the metric $`{\displaystyle \frac{g_{ii}^{KdV}(𝜷)}{2(\beta ^i)^2}}`$ of constant curvature $`R_{ij}^{ij}={\displaystyle \frac{1}{2}}`$. Let $`J_4`$ be the nonlocal Hamiltonian operator of the form (36) associated to the conformally flat metric $`{\displaystyle \frac{g_{ii}^{KdV}(𝜷)}{(\beta ^i)^3}}`$ with Riemannian curvature $`R_{ij}^{ij}={\displaystyle \frac{1}{8}}(w_+^i+w_+^j),`$ where $$w_+^i=\left(\beta ^1+\beta ^2+\beta ^3+\frac{2_{ji}(\beta ^i\beta ^j)}{\beta ^i+\alpha _1}\right),$$ (57) with $`\alpha _1`$ defined in (52). The $`w_+^i`$, $`i=1,\mathrm{},3`$, are the velocities (originally obtained by Whitham ) of the usual positive KdV modulated flow $$\frac{\beta ^i}{\tau _+}+w_+^i(𝜷)\frac{\beta ^i}{y}=0.$$ ###### Remark 4.4 The Hamiltonian operator $`J_1`$ corresponds to the average over a one-dimensional torus of the Gardner-Zakharov KdV Hamiltonian structure $`𝒫_1=8_y`$ while the Hamiltonian operator $`J_2`$ corresponds to the average over a one-dimensional torus of the Lenard-Magri local Hamiltonian structure $`𝒫_2=2_{yyy}+2U_y+2_yU`$. We use this unusual normalization in order to be consistent with the normalization of the CH equation. Defining the recursive operator $`=𝒫_2(𝒫_1)^1`$, we obtain the family of non-local Hamiltonian operators $`𝒫_{k+1}=𝒫_k,`$ $`k1`$. The averaged KdV non-local Hamiltonian operators $`J_3`$ and $`J_4`$ corresponds to the average over a one-dimensional torus of the non-local Hamiltonian operators $`𝒫_3`$ and $`𝒫_4`$, respectively. The modulation equations of the negative KdV flow can be written in Hamiltonian form with Hamiltonian operator $`J_1`$ and Hamiltonian density $`_0`$, which is the average over the one dimensional torus of the Casimir generating the KdV negative flow, namely $$_0=\frac{d\mathrm{\Theta }}{\rho (\mathrm{\Theta })}.$$ (58) ###### Lemma 4.5 In the $`\beta `$s’ coordinates $$_0=ip(0),$$ where $`p(\eta )`$ is the Abelian integral of the quasi-momentum $`dp(\eta )`$ defined in (52). To prove the lemma, we compute the integral in (58) in the $`\beta `$s’ coordinates obtaining $$_0=\frac{d\mathrm{\Theta }}{\rho (\mathrm{\Theta })}=\sqrt{\beta ^1\beta ^2\beta ^3}\alpha _0=i\frac{(\mathrm{\Lambda }_0(\eta ))}{2d\xi }\|_{\xi =0},\eta =\frac{1}{\xi ^2},$$ where $$\mathrm{\Lambda }_0(\eta )=\sqrt{\beta ^1\beta ^2\beta ^3}\frac{{\displaystyle \frac{1}{\eta }}+\alpha _0}{\sqrt{(\eta \beta ^1)(\eta \beta ^2)(\eta \beta ^3)}}d\eta .$$ (59) $`\mathrm{\Lambda }_0(\eta )`$ is a normalized third kind differential with first order poles at the points $`O^\pm =(0,\pm \sqrt{\beta ^1\beta ^2\beta ^3})`$ with residue $`\pm 1`$, respectively. The constant $`\alpha _0`$ in (59) is uniquely determined by the condition $`{\displaystyle _a}\mathrm{\Lambda }_0(\eta )=0.`$ From the Riemann bilinear relation, it is finally immediate to verify that $$_0=i\frac{(\mathrm{\Lambda }_0(\eta ))}{2d\xi }\|_{\xi =0}=ip(0),\eta =\frac{1}{\xi ^2},$$ (60) where $`p(\eta )`$ is the Abelian integral of the quasi-momentum $`dp(\eta )`$ defined in (52). Finally the following result holds. ###### Lemma 4.6 The first negative KdV averaged flow (53) can be written in the Hamiltonian form $$\beta _\tau _{}^i=J_1^{ij}\frac{\delta _0}{\delta \beta ^j}=J_2^{ij}\frac{\delta _1}{\delta \beta ^j}=J_3^{ij}\frac{\delta _2}{\delta \beta ^j}=J_4^{ij}\frac{\delta _3}{\delta \beta ^j},$$ (61) where $`_0`$ is the Casimir of $`J_2`$ defined in (58) and the Hamiltonian densities are determined by the recursion scheme $$J_1^{ij}\frac{\delta _s}{\delta \beta ^j}=J_2^{ij}\frac{\delta _{s1}}{\delta \beta ^j},s0.$$ The Hamiltonian $`_s`$, $`s1`$, are generated by the expansion for $`\eta 0`$ of the quasi-momentum $`dp(\eta )`$. Indeed, $$i_{(\eta ,w)}^{(\eta ,w)}𝑑p(\xi )=\frac{1}{2}(_0+\eta _1+\eta ^2_2+\frac{3}{4}\eta ^3_3+\mathrm{}),\eta 0.$$ (62) ### 4.3 CH versus KdV modulation equations In the previous section we computed the modulation equations of the first negative KdV flow. In this subsection, we compute the average of the reciprocal transformation defined in section 4.1 and we show that the negative KdV modulation equations are transformed to the CH modulation equations. Finally, both KdV-Whitham and CH-Whitham systems are Hamiltonian systems, so we end the section investigating how the reciprocal transformation acts on the Hamiltonian structures of the two systems. To compare KdV and CH, we first need to reduce the even spectral curve of CH to the odd spectral curve of KdV. The natural change of coordinates $`(\lambda ,y)(\eta ,w)`$ $$\eta =\frac{1}{\lambda +\nu },w^2=\frac{y^2}{(\lambda +\nu )^4_{i=1}^3(\nu +u^i)},$$ maps the even spectral curve $`y^2=(\lambda +\nu )(\lambda u^1)(\lambda u^2)(\lambda u^3)`$ to the odd KdV spectral curve $$w^2=(\eta \beta ^i)(\eta \beta ^2)(\eta \beta ^3),\beta ^i=\frac{1}{\nu +u^i}.$$ The differentials $`\mathrm{\Omega }_\nu (\lambda )`$ and $`\sigma _1(\lambda )`$ defined in (24) and (15) transform to $$\mathrm{\Omega }_\nu (\lambda )\frac{1}{2}dp(\eta ),$$ (63) $$\sigma _1(\lambda )\mathrm{\Lambda }_0(\eta ),$$ (64) with $`dp(\eta )`$ as in (52) and $`\mathrm{\Lambda }_0(\eta )`$ as in (59). It follows from (63) and (64) that the change of coordinates $`\beta ^i=1/(\nu +u^i)`$ transforms the speeds $`C^i(𝒖)`$ defined in (25) to $$\stackrel{~}{C}^i(𝜷)=\frac{1}{\beta ^1}+\frac{1}{\beta ^2}+\frac{1}{\beta ^3}\nu +2\frac{\alpha _0_{ji,j=1}^3(\beta ^i\beta ^j)}{\beta ^i(\beta ^i+\alpha _1)}.$$ (65) Now we show that the averaged reciprocal transformation maps the modulation equations of the KdV negative flow to the CH modulation equations. Indeed averaging over a period the inverse of (43) $$dx=1/\rho dy+𝚞d\tau _{},dt=d\tau _{},$$ we get the averaged reciprocal transformation $$dx=dy\frac{d\mathrm{\Theta }}{\rho (\mathrm{\Theta })}+d\tau _{}𝚞(\mathrm{\Theta })𝑑\mathrm{\Theta }dt=d\tau _{}.$$ (66) ###### Proposition 4.7 The averaged reciprocal transformation (66) takes the form $$dx=_0dy+𝒩d\tau _{},dt=d\tau _{},$$ (67) where $`_0`$ is the Casimir defined in (58) and $$𝒩=\frac{1}{\beta ^1}+\frac{1}{\beta ^2}+\frac{1}{\beta ^3}\nu +2\alpha _0=\frac{1}{2}(_0)^2\nu ,$$ (68) where $`(_0)^2=_i(g_{ii}^{KdV})^1(_{\beta ^i}_0)^2`$. Finally (67) is a reciprocal transformation for the KdV-Whitham negative flow. To prove the proposition we observed that the one form (67) is closed by (68). The proof of (68) follows from the identity $`_0=ip(0)`$, and the variational formula $`_{\beta ^i}p(0)={\displaystyle \frac{1}{4}}dp(\beta ^i)\mathrm{\Lambda }_0(\beta ^i)`$. Next we show that the modulation equations of the first negative KdV flow are mapped by the reciprocal transformation (67) to the CH modulation equations. ###### Proposition 4.8 The reciprocal transformation $`dx=_0dy+𝒩d\tau _{}`$, $`dt=d\tau _{}`$, where $`_0`$ and $`𝒩`$ are as in (58) and (68) respectively, transforms the modulation equations (53) of the first negative KdV flow $$_\tau _{}\beta ^i+v^i(𝜷)_y\beta ^i=0,$$ (69) where $`v^i(𝛃)`$, $`i=1,\mathrm{},3`$ are as in (54) to the CH modulation equations $$_t\beta ^i+\stackrel{~}{C}^i(𝜷)_x\beta ^i=0,$$ (70) where the CH velocities are defined in (65). Viceversa, the inverse reciprocal transformation $`dy={\displaystyle \frac{1}{_0}}dx{\displaystyle \frac{𝒩}{_0}}dt`$ transforms (70) into (69). Indeed plugging (67) into (69) we get $`_t\beta ^i+V^i(𝜷)_x\beta ^i=0`$, where $$V^i(𝜷)=v^i(𝜷)_0+𝒩=\stackrel{~}{C}^i(𝜷).$$ To compare the CH and KdV Hamiltonian structures, we express the CH-Whitham Hamiltonian operators introduced in subsection 3.1 in the $`\beta `$s’ coordinates. In these coordinates, using (60), (62) and (63), the CH flat metric $`g_{ii}(𝒖)`$ defined in (27) takes the form $$\stackrel{~}{g}_{ii}(𝜷)=\frac{1}{(\beta ^i)^3}\frac{\underset{[}{Res}\eta =\beta ^i]\{{\displaystyle \frac{(dp(\eta ))^2}{d\eta }}\}}{p(0)^2/4}=\frac{g_{ii}^{KdV}(𝜷)}{(\beta ^i)^3_0^2},i=1,2,3.$$ (71) The other metrics defined in (28) transform, respectively, to $$\frac{g_{ii}^{KdV}(𝜷)}{2(\beta ^i)^2_0^2},\frac{g_{ii}^{KdV}(𝜷)}{4\beta ^i_0^2},\frac{g_{ii}^{KdV}(𝜷)}{8_0^2},$$ (72) and are, respectively, flat, of constant curvature and conformally flat. Let $`\stackrel{~}{A}_1`$ and $`\stackrel{~}{A}_2`$ be the local Hamiltonian operators associated to the CH flat metrics $`\stackrel{~}{g}_{ii}(𝜷)`$ and $`{\displaystyle \frac{1}{2}}\stackrel{~}{g}_{ii}(𝜷)\beta ^i`$, respectively. Moreover, let $`\stackrel{~}{A}_3`$ and $`\stackrel{~}{A}_4`$ be the non-local Hamiltonian operators associated to the CH constant curvature and conformally flat metrics, $`{\displaystyle \frac{1}{4}}\stackrel{~}{g}_{ii}(𝜷)(\beta ^i)^2`$ and $`{\displaystyle \frac{1}{8}}\stackrel{~}{g}_{ii}(𝜷)(\beta ^i)^3`$, respectively. Then, in the $`\beta `$s’ coordinates the CH modulation equations in Hamiltonian form are $$\beta _t^i=\stackrel{~}{A}_1^{ij}\frac{\delta \stackrel{~}{h}_2}{\delta \beta ^j}=\stackrel{~}{A}_2^{ij}\frac{\delta \stackrel{~}{h}_1}{\delta \beta ^j}=\stackrel{~}{A}_3^{ij}\frac{\delta \stackrel{~}{h}_0}{\delta \beta ^j}=\stackrel{~}{A}_4^{ij}\frac{\delta \stackrel{~}{h}_1}{\delta \beta ^j}$$ (73) where the Hamiltonians $`\stackrel{~}{h}_j`$, $`i=1,\mathrm{},2,`$ are the averaged conservation laws introduced in subsection 3.1 expressed in the $`\beta `$s’ coordinates. Indeed, following (62), the positive CH Hamiltonians are obtained from the coefficients of the expansion for $`\eta 0`$ of the differential $`{\displaystyle \frac{dp(\eta )}{_0}}`$ and take the form $$\stackrel{~}{h}_1=1\frac{\nu }{_0},\stackrel{~}{h}_0=\frac{_1}{_0}\nu ,\stackrel{~}{h}_j=\frac{_{j1}}{_0}\nu \frac{_j}{_0},j1,$$ (74) where the $`_k`$ are defined in (62). The following theorem by Ferapontov and Pavlov describes the action of a reciprocal transformation for an Hamiltonian hydrodynamic equation and we use it to clarify the relation between the Hamiltonian structures of the averaged KdV and CH Hamiltonian structures. ###### Theorem 4.9 Let $`\beta _\tau ^i=J^{ij}{\displaystyle \frac{h}{\beta ^j}}`$ be an Hamiltonian system associated to the local operator $`J^{ij}=g^{ii}\delta _i^j{\displaystyle \frac{d}{dy}}g^{ii}\mathrm{\Gamma }_{ik}^j\beta _y^k`$. Then, under the action of the reciprocal transformation $`dx=Ady+Bd\tau ,dt=d\tau `$, where $`d(Ady+Bd\tau )=0`$, the transformed system is Hamiltonian $`\beta _t^i=\stackrel{~}{J}^{ij}{\displaystyle \frac{\stackrel{~}{h}}{\beta ^j}},`$ with nonlocal operator $$\stackrel{~}{J}^{ij}=\stackrel{~}{g}^{ii}\delta ^{ij}\frac{d}{dx}\stackrel{~}{g}^{ii}\stackrel{~}{\mathrm{\Gamma }}_{ik}^j\beta _x^k+\stackrel{~}{w}^i\beta _x^i\left(\frac{d}{dx}\right)^1\beta _x^j+\beta _x^i\left(\frac{d}{dx}\right)^1\stackrel{~}{w}^j\beta _x^j,$$ and hamiltonian density $`\stackrel{~}{h}=h/A`$. The transformed metric is $`\stackrel{~}{g}_{ii}=g_{ii}/A^2`$, $`\stackrel{~}{\mathrm{\Gamma }}`$ is the Levi–Civita connection and $$\stackrel{~}{w}^i=^i_iAA\frac{1}{2}(A)^2=g^{ii}\left(_i^2A\underset{j}{}\mathrm{\Gamma }_{ii}^j_jA\right)A\frac{1}{2}\underset{j}{}g^{jj}(_jA)^2.$$ (75) Moreover, the transformed metric is conformally flat with curvature tensor $$\stackrel{~}{R}_{ij}^{ij}=\stackrel{~}{w}^i+\stackrel{~}{w}^j.$$ By the above theorem and by inspection of (61) and (73), where the metrics have been defined in (71), (72) and (56), we conclude that the local KdV-Whitham Poisson operators, $`J_1`$ and $`J_2`$, are mapped to the nonlocal CH-Whitham Poisson operators $`\stackrel{~}{A}_4`$ and $`\stackrel{~}{A}_3`$, respectively, by the reciprocal transformation $`dx=_0dy+𝒩d\tau _{}`$, $`dt=d\tau _{}`$. ###### Remark 4.10 The application of theorem 4.9 deserves a special attention in the case $`\nu 0`$ for the computation of the Hamiltonian densities $`\stackrel{~}{h}_k`$. Indeed neither the KdV–Hamiltonian operators $`J_i`$ nor the averaged negative KdV Hamiltonian densities $`_s`$, $`s=0,1,2,3`$, depend on the parameter $`\nu `$ while the CH averaged Hamiltonian densities $`\stackrel{~}{h}_s`$, $`s=1,0,1,2,`$ in (74) do. To solve this apparent contradiction, we recall that the Hamiltonian densities are defined modulo Casimirs. In particular, the Casimir $`\nu _0`$ of the second KdV-Whitham Hamiltonian operator is transformed to the CH Hamiltonian density $`\nu `$ which generates the constant flow term associated to $`\stackrel{~}{A}_3`$. Similarly, the term $`\nu `$ is mapped by the reciprocal transformation to the CH Hamiltonian density $`\nu /_0`$ which generates the constant flow term associated to $`\stackrel{~}{A}_4`$. Applying theorem 4.9 to the local CH–Whitham Poisson operators, we prove that the local Hamiltonian CH operators $`\stackrel{~}{A}_1`$ and $`\stackrel{~}{A}_2`$, are mapped to the nonlocal KdV–Whitham Poisson operators $`J_4`$ and $`J_3`$, respectively, by the reciprocal transformation $`dy=1/_0dx𝒩/_0dt`$. By the same theorem 4.9, the corresponding KdV Hamiltonian densities are related to the CH ones by $$\stackrel{~}{h}_2_0=_3,\stackrel{~}{h}_1_0=_2,$$ where the above identities again hold modulo Casimirs of the corresponding KdV Hamiltonian operators. We summarize the results of this section in the following. ###### Theorem 4.11 The reciprocal transformation $`dx=_0dy+𝒩d\tau _{}`$, $`dt=d\tau _{}`$, maps the KdV local Hamiltonian operators $`J_1`$, $`J_2`$ and the corresponding Hamiltonian densities $`_0`$, $`_1`$ to $$J_1^{ik}\frac{\delta _0}{\delta \beta ^k}\stackrel{~}{A}_4^{ik}\frac{\delta \stackrel{~}{h}_1}{\delta \beta ^k}$$ $$J_2^{ik}\frac{\delta _1}{\delta \beta ^k}\stackrel{~}{A}_3^{ik}\frac{\delta \stackrel{~}{h}_0}{\delta \beta ^k}$$ where $`\stackrel{~}{A}_4`$ is the CH nonlocal Hamiltonian operator associated to the covariant conformally flat metric $`{\displaystyle \frac{g_{ii}^{KdV}(𝛃)}{8_0^2}}`$ and $`\stackrel{~}{h}_1=1\nu /_0`$ the corresponding CH Hamiltonian density ( $`\stackrel{~}{h}_1_0=_0`$ modulo Casimir of $`J_1`$). The Hamiltonian operator $`\stackrel{~}{A}_3`$ is the CH nonlocal Hamiltonian operator associated to the constant curvature metric $`{\displaystyle \frac{g_{ii}^{KdV}(𝛃)}{4_0^2\beta ^i}}`$ and $`\stackrel{~}{h}_0=\nu +{\displaystyle \frac{_1}{_0}}`$ the corresponding CH Hamiltonian density ( $`\stackrel{~}{h}_0_0=_1`$ modulo a Casimir of $`J_2`$). The reciprocal transformation $`dy=1/_0dx𝒩/_0dt`$, $`d\tau _{}=dt`$ maps the local CH Hamiltonian operators $`\stackrel{~}{A}_1`$, $`\stackrel{~}{A}_2`$, and corresponding Hamiltonian densities $`\stackrel{~}{h}_2={\displaystyle \frac{_3}{_0}}\nu {\displaystyle \frac{_2}{_0}}`$, $`\stackrel{~}{h}_1={\displaystyle \frac{_2}{_0}}\nu {\displaystyle \frac{_1}{_0}}`$ to $$\stackrel{~}{A}_1^{ik}\frac{\delta \stackrel{~}{h}_2}{\delta \beta ^k}J_4^{ik}\frac{\delta _3}{\delta \beta ^k}$$ $$\stackrel{~}{A}_2^{ik}\frac{\delta \stackrel{~}{h}_1}{\delta \beta ^k}J_3^{ik}\frac{\delta _2}{\delta \beta ^k}$$ where the nonlocal KdV Hamiltonian operator $`J_4`$ is associated to the covariant conformally flat metric $`{\displaystyle \frac{g_{ii}^{KdV}(𝜷)}{(\beta ^i)^3}}`$ and $`\stackrel{~}{h}_2_0=_3`$ modulo a Casimir of $`J_4`$. The nonlocal KdV Hamiltonian operator $`J_3`$ is associated to the covariant constant curvature metric $`{\displaystyle \frac{g_{ii}^{KdV}(𝜷)}{2(\beta ^i)^2}}`$ and $`\stackrel{~}{h}_1_0=_2`$ modulo a Casimir of $`J_3`$. We conclude this section by illustrating with a table all the reciprocal transformations between the various Hamiltonian operators. Acknoledgments We express our gratitude to A. Maltsev for the many useful discussions and for pointing out to us the Lagrangian averaging method for the Camassa-Holm equations. We also thank B. Dubrovin, G. Falqui and D. Holm for the useful discussions and comments of the manuscript. This work has been partially supported by the GNFM-INdAM Project ”Onde nonlineari, struttura tau e geometria delle varietà invarianti: il caso della gerarchia di Camassa-Holm” and by the European Science Foundation Programme MISGAM (Method of Integrable Systems, Geometry and Applied Mathematics).
warning/0506/nucl-ex0506005.html	ar5iv	text	# 𝑄²-Dependence of the Neutron Spin Structure Function 𝑔₂^𝑛 at Low 𝑄² ## Abstract We present the first measurement of the $`Q^2`$-dependence of the neutron spin structure function $`g_2^n`$ at five kinematic points covering $`0.57`$ $`(\mathrm{GeV}/c)^2`$ $`Q^21.34`$ $`(\mathrm{GeV}/c)^2`$ at $`x0.2`$. Though the naive quark-parton model predicts $`g_2=0`$, non-zero values occur in more realistic models of the nucleon which include quark-gluon correlations, finite quark masses or orbital angular momentum. When scattering from a non-interacting quark, $`g_2^n`$ can be predicted using next-to-leading order fits to world data for $`g_1^n`$. Deviations from this prediction provide an opportunity to examine QCD dynamics in nucleon structure. Our results show a positive deviation from this prediction at lower $`Q^2`$, indicating that contributions such as quark-gluon interactions may be important. Precision data obtained for $`g_1^n`$ are consistent with next-to-leading order fits to world data. Over the past 30 years, significant progress has been made in understanding the spin structure of the nucleon through measurements using polarized deep-inelastic lepton scattering (DIS). Most of these experiments were focused on precise measurements of the spin structure function $`g_1`$. In the naive quark-parton model (QPM), $`g_1`$ is directly related to the contributions of the individual quark flavors to the overall spin of the nucleon (see e.g. Ref. Thomas ). Sum rules based on this simple model have provided fertile ground for understanding the origin of the nucleon spin in terms of quark degrees of freedom. In addition, next-to-leading-order (NLO) analyses of the world $`g_1`$ data (see e.g. Refs. BB ; AAC03 ) have provided indirect information about the role of gluons in the nucleon’s spin. The QPM is expected to be valid in the scaling limit, where the four-momentum transfer squared $`Q^2`$ and energy transferred $`\nu `$ approach infinity. As $`Q^2`$ becomes large, the electron-nucleon interaction can be described by scattering from a massless, non-interacting quark carrying a fraction $`x=Q^2/2M\nu `$ of the nucleon’s momentum ($`M`$ is the nucleon mass). At finite $`Q^2`$ and $`\nu `$, effects such as gluon Bremsstrahlung, vacuum polarization and vertex corrections can be accurately calculated using perturbative QCD. For $`g_1`$, non-perturbative QCD processes, such as quark-gluon and quark-quark correlations, are suppressed relative to the asymptotically-free contributions by factors of $`1/Q`$ and $`1/Q^2`$, respectively. Polarized DIS also provides information about a second spin structure function, $`g_2`$, which is identically zero in the QPM Manohar . Interest in $`g_2`$ arises because, unlike $`g_1`$, contributions from certain non-perturbative processes enter at the same order in $`Q^2`$ as the asymptotically-free contributions. An appropriate formalism for understanding $`g_2`$ is the operator product expansion (OPE) OPE1 ; OPE2 which is a model-independent approach based directly on QCD. Here, the unknown hadronic currents relevant for polarized DIS are expanded in terms of quark and gluon operators and grouped by factors of $`(1/Q)^{\tau 2}`$, where $`\tau =2,3,4,\mathrm{}`$ is known as the twist of the operator. Both $`g_1`$ and $`g_2`$ contain a contribution from a twist-2 operator that corresponds to scattering from a massless, non-interacting quark. Operators with higher-twist represent contributions from non-perturbative processes such as quark-quark and quark-gluon correlations, and from quark mass effects. Because these correlations are responsible for quark confinement, higher-twist effects must be included in any realistic model of the nucleon OPE2 . Using the fact that $`g_1`$ and $`g_2`$ contain the same twist-2 operator, Wandzura and Wilczek g2ww derived the following expression for the asymptotically-free contribution to $`g_2`$, in terms of $`g_1`$, $$g_2^{WW}(x,Q^2)=g_1(x,Q^2)+_x^1\frac{g_1(x,Q^2)}{x}𝑑x.$$ (1) The world data for $`g_1`$ cover a broad range in $`x`$ and $`Q^2`$ with relatively high precision. Polarized parton distributions can be extracted from NLO fits to the data and evolved to different values of $`Q^2`$ using the Dokshitzer-Gribov-Lipatov-Alarelli-Paresi (DGLAP) procedure DGLAP1 ; DGLAP2 ; DGLAP3 ; DGLAP4 . Because the world data are at relatively large $`Q^2`$, where higher-twist effects should be negligible, these evolved parton distributions allow one to calculate the twist-2 part of $`g_1`$, and therefore $`g_2^{WW}`$, in most kinematic regions accessible today. Precise measurements of $`g_2`$ at specific values of $`x`$ and $`Q^2`$ can be compared to $`g_2^{WW}`$, providing a unique opportunity to cleanly isolate higher-twist contributions. Previous measurements of $`g_2`$ E143 ; E154 ; E155 ; E155x were aimed at testing OPE sum rules and covered a wide range in $`x`$, at an average $`Q^2`$ of 5 $`(\mathrm{GeV}/c)^2`$. Proton data show general agreement with $`g_2^{WW}`$ but lack the precision needed to make a definitive statement about the size of higher-twist effects by direct comparison. Neutron results have much larger uncertainties and cannot distinguish between $`g_2^{WW}`$ and $`g_2=0`$. We present a new measurement of the $`Q^2`$-dependence of $`g_2`$ for the neutron at low $`Q^2`$, while keeping $`x`$ approximately constant. With statistical uncertainties more than an order of magnitude smaller than existing data, these results allow, for the first time, a precise comparison with $`g_2^{WW}`$ to study the $`Q^2`$-dependence of higher-twist effects. Longitudinally-polarized electrons were scattered from a polarized <sup>3</sup>He target in Hall A HALLA at the Thomas Jefferson National Accelerator Facility (Jefferson Lab) through the inclusive process $`\stackrel{}{{}_{}{}^{3}\mathrm{He}}(\stackrel{}{e},e^{})`$. We measured $`g_2`$ at five values of $`Q^2`$ between $`0.57`$ $`(\mathrm{GeV}/c)^2`$ and $`1.34`$ $`(\mathrm{GeV}/c)^2`$, with $`x0.2`$. The kinematics for this $`x`$ are well-matched to the beam energies available at Jefferson Lab and are in a region where $`g_2^{WW}`$ is relatively large. Low $`Q^2`$ provides access to the region where higher-twist effects are expected to become important. The invariant mass squared of the photon-nucleon system, $`W^2=M^2+2M\nu Q^2`$, was kept $`>3.8`$ $`(\mathrm{GeV}/c)^2`$ to minimize resonance contributions. The spin of the <sup>3</sup>He nuclei could be oriented either longitudinal or transverse to the incoming electron spin. The parallel and perpendicular cross-section differences, $`\mathrm{\Delta }\sigma _{}`$ and $`\mathrm{\Delta }\sigma _{}`$, are defined in terms of helicity-dependent cross-sections $`\sigma ^{ij}`$ as, $$\mathrm{\Delta }\sigma _{}=\sigma ^{}\sigma ^{},\mathrm{\Delta }\sigma _{}=\sigma ^{}\sigma ^{},$$ (2) where the single and double arrows refer to the beam and target spin directions, respectively, relative to the incident electron momentum direction. An up (down) arrow refers to spin aligned parallel (anti-parallel) to the momentum. The double right arrow refers to target polarization in the scattering plane, perpendicular to the incident electron momentum, with the scattered electron detected on the side of the beam line towards which the target spin points. These cross-section differences are related to $`g_1`$ and $`g_2`$ by Thomas , $`\mathrm{\Delta }\sigma _{}`$ $`=`$ $`{\displaystyle \frac{4\alpha ^2E^{}}{Q^2EM\nu }}[(E+E^{}\mathrm{cos}\theta )g_1(x,Q^2)`$ (3) $``$ $`2xMg_2(x,Q^2)],`$ $`\mathrm{\Delta }\sigma _{}`$ $`=`$ $`{\displaystyle \frac{4\alpha ^2E^{}}{Q^2EM\nu }}E^{}\mathrm{sin}\theta [g_1(x,Q^2)`$ (4) $`+`$ $`{\displaystyle \frac{4xEM}{Q^2}}g_2(x,Q^2)],`$ where $`E`$ and $`E^{}`$ are the incident and scattered electron energies and $`\theta `$ is the laboratory scattering angle. Polarized electrons were produced by photo-emission from a strained GaAs crystal using circularly polarized laser light. The average beam polarization was $`P_b=(76\pm 3)\%`$ as measured using both Møller and Compton polarimeters HALLA and the average beam current was $`12.0`$ $`\mu `$A, with an uncertainty of $`\pm 1\%`$. The helicity of the electron beam was flipped on a pseudo-random basis at a rate of 30 Hz to minimize helicity-correlated systematic effects. The beam energy ranged from $`3.5`$ GeV to $`5.7`$ GeV and was measured with a relative uncertainty below $`10^3`$ using both elastic electron-proton scattering and a magnetic field measurement. Scattered particles were detected by either one of a pair of magnetic spectrometers arranged symmetrically on either side of the beam line HALLA . The momentum of the particles was determined by reconstructing their trajectories using drift chambers. Electrons were identified using a threshold gas Cherenkov detector and a two-layer lead-glass electromagnetic calorimeter. The largest background was from $`\pi ^{}`$ photo-production in the target. The ratio of pion rate to electron rate detected in the spectrometer was $`<3.3`$. The ratio of the pion asymmetry to the electron asymmetry was $`<4.1`$. The pion rejection factor was $`>10^4:1`$ with an electron detection efficiency of $`97\%`$, which was sufficient to reduce the pion contamination to a negligible level. Electrons scattered in the entrance and exit windows of the target cell were removed using software cuts. To study polarized neutrons, a target containing <sup>3</sup>He nuclei was polarized using spin-exchange optical pumping SEOP . The <sup>3</sup>He ground state is dominated by the $`S`$-state in which the two proton spins are anti-aligned and the spin of the nucleus is carried entirely by the neutron. The <sup>3</sup>He is contained in a sealed, two-chambered, aluminosilicate glass cell, along with a small quantity of N<sub>2</sub> and Rb to aid in the polarization process. Polarized <sup>3</sup>He is produced in the spherical upper chamber by first polarizing Rb atoms with optical pumping. These atoms can transfer their spin to the <sup>3</sup>He nucleus during binary collisions. Incident electrons scatter from the polarized <sup>3</sup>He in the cylindrical lower chamber which is 40 cm long with side walls of thickness $`1.0`$ mm and end windows of thickness $`130`$ $`\mu `$m. The <sup>3</sup>He density as seen by the beam is $`2.9\times 10^{20}`$/cm<sup>3</sup>. The average in-beam target polarization was $`P_t=(40.0\pm 1.4)\%`$ as measured using both nuclear magnetic resonance NMR and electron paramagnetic resonance EPR . The <sup>3</sup>He spin structure functions $`g_1`$ and $`g_2`$ were obtained using our measured cross-section differences and Eqs. (3) and (4). The ratios of these cross-section differences to the total unpolarized cross-section are defined as the longitudinal and transverse asymmetries, respectively. To achieve our desired statistical precision on $`g_2^n`$, the raw (uncorrected) transverse asymmetries were measured to a statistical uncertainty of $`10^4`$. It was also necessary to keep sources of false asymmetries at, or below, this level. A feedback system was used to keep helicity-dependent beam charge asymmetry below $`50\times 10^6`$ for a typical run. Quasi-elastic polarized electron scattering from an unpolarized <sup>12</sup>C target was used to measure the false asymmetry and yielded $`(67\pm 46)\times 10^6`$. Misalignment of the target polarization direction was typically less than $`\pm 0.3^{}`$, implying a negligible contribution from $`\mathrm{\Delta }\sigma _{}`$ when measuring $`\mathrm{\Delta }\sigma _{}`$. For each kinematic setting, equal quantities of data were collected with the overall sign of the electron helicities reversed and/or the target polarization rotated by $`180^{}`$. The absolute value of the measured asymmetries for each of the beam helicity and target spin combinations were consistent with each other within the statistical uncertainties, indicating that there were no significant false asymmetries. The raw cross-section differences were determined with a relative uncertainty of $`6\%`$ and were checked against elastic data taken during this experiment. Spectrometer acceptances were provided by the Hall A Single-Arm Monte Carlo simulation Deur ; Kramer used for previous polarized <sup>3</sup>He experiments E94010 ; A1nref . Corrections were applied for beam and target polarizations, $`P_b`$ and $`P_t`$, and for radiative effects using a modified version of the code POLRAD 2.0 POLRAD for internal corrections, and the formalism in Refs. MoTsai ; Stein for external corrections. Input to the radiative corrections came from fits to measured polarized cross-sections in the quasi-elastic region Slifer , fits to measured $`g_1^n`$ and $`g_2^n`$ data in the resonance region E94010 and a fit to $`g_1^n/F_1^n`$ A1nref , along with calculations of $`g_2^{WW}`$, in the deep-inelastic region. Our measured data were used to guide the fits near $`x0.2`$, and the systematic uncertainty in the radiative corrections is dominated by uncertainties in the fits where data are sparse. To obtain neutron spin structure functions, a correction was applied to the measured values of $`g_1^{{}_{}{}^{3}\mathrm{He}}`$ and $`g_2^{{}_{}{}^{3}\mathrm{He}}`$ based on a model <sup>3</sup>He wavefunction Bissey , $$g_{1,2}^n=\frac{1}{P_n+0.056}\left[g_{1,2}^{{}_{}{}^{3}\mathrm{He}}+\left(0.0142P_p\right)g_{1,2}^p\right],$$ (5) where $`P_p=0.028_{0.020}^{+0.036}`$ and $`P_n=0.86_{0.004}^{+0.009}`$ are the effective proton and neutron polarizations in <sup>3</sup>He Bissey ; Ciofi ; Nogga . For calculations of $`g_1^p`$ we used the average of scenarios 1 and 2 in the Blümlein and Böttcher (BB) NLO fit to world data BB , evolved to our $`Q^2`$ using the DGLAP procedure DGLAP1 ; DGLAP2 ; DGLAP3 ; DGLAP4 . For $`g_2^p`$, values for $`g_2^{WW}`$ were calculated using BB and Eq. (1). Results for $`g_1^n`$ and $`g_2^n`$ from this experiment are given in Table 1. The largest contributions to the systematic uncertainties in $`g_1^n`$ and $`g_2^n`$ come from the uncertainties in the cross-section differences, and, for $`g_2^n`$, from uncertainties in the models for the transversely polarized quasi-elastic and resonance-region tails used for radiative corrections. Figure 1 shows our data for $`g_1^n`$ as a function of $`Q^2`$. Also shown are $`g_1^n`$ predictions from the NLO analyses of BB and M. Hairi et al.(AAC03e) AAC03 . The uncertainty in the BB curve comes from propagating the uncertainties in the parton distribution functions used in the fit. The world data for $`g_1^n`$ at our $`x`$ are at significantly larger $`Q^2`$ where higher-twist contributions should be suppressed. The good agreement with these evolved fits indicates that, within the uncertainties, we are not seeing higher-twist effects in $`g_1^n`$. This also gives confidence that we can use the BB fit to calculate $`g_2^{WW}`$ without introducing significant higher-twist effects from $`g_1`$. Figure 2 shows our data for $`g_2^n`$ as a function of $`Q^2`$ along with calculations of $`g_2^{WW}`$ from the BB and AAC03e NLO analyses. The data are more than $`5\sigma `$ above zero, and at lower $`Q^2`$, show a systematic positive deviation from $`g_2^{WW}`$, which we interpret as evidence for non-zero higher-twist contributions. In the OPE, the twist-2 and twist-3 contributions to $`g_2`$ enter at the same order in $`Q^2`$, with additional higher-twist contributions suppressed by powers of $`1/Q`$. Assuming these additional higher-twist contributions are small, we expect the quantity $`g_2g_2^{WW}`$ to be constant as a function of $`Q^2`$. When compared to $`g_2^{WW}`$ from BB, a fit to our data gives $`g_2g_2^{WW}=0.0262\pm 0.0043\pm 0.0080\pm 0.0099`$, with a reduced $`\chi ^2`$ of 1.4. The first two uncertainties are from the experiental statistical and systematic uncertainties and the third is from the uncertainty in the BB fit. Attempts to fit the data with functions that contain additional higher-twist dependence and/or logarithmic perturbative QCD corrections did not improve the quality of the fits. Also shown in Figure 2 are chiral soliton model calculations from Weigel et al. Weigel1 ; Weigel2 and Wakamatsu Wakamatsu , and a bag model calculation of the higher-twist contribution from Stratmann Strat , combined with $`g_2^{WW}`$ from BB. Our data indicate a positive contribution from higher-twist effects while the model calculations generally predict a negative contribution. Figure 3 shows the improved quality of our $`g_2^n`$ data (solid points) plotted versus $`x`$, as compared to older data (open squares) from the Stanford Linear Accelerator Center E155x . Also shown are recent data (open triangles) from Jefferson Lab at higher $`x`$ A1nref . The data shown cover a range from $`Q^2=0.7`$ $`(\mathrm{GeV}/c)^2`$ at the lowest $`x`$ to $`Q^2=20`$ $`(\mathrm{GeV}/c)^2`$ at the highest $`x`$. For reference, $`g_2^{WW}`$ from BB is also included at constant $`Q^2=1.0`$ $`(\mathrm{GeV}/c)^2`$. In summary, we have made the first measurement of the $`Q^2`$-dependence of the $`g_2`$ spin structure function for the neutron at $`x0.2`$ in the range $`Q^2=0.571.34`$ $`(\mathrm{GeV}/c)^2`$. With a factor of $`>10`$ improvement in statistical precision over previous measurements, these data allow for an accurate determination of the higher-twist contributions by direct comparison with the twist-2 $`g_2^{WW}`$ prediction. Our results show a positive deviation from $`g_2^{WW}`$ at lower-$`Q^2`$, indicating a non-zero higher-twist contribution, and are inconsistent with model calculations, which generally predict a negative higher-twist contribution. Precision data for $`g_1`$ obtained in this kinematic range showed good agreement with the NLO analyses of world data, indicating no significant higher-twist contributions within the uncertainties. ###### Acknowledgements. We thank the personnel of Jefferson Lab for their support and efforts towards the successful completion of this experiment. We would also like to thank M. Stratmann, H. Weigel, L. Gamberg and M. Wakamatsu for additional model calculations at our kinematics. This work was supported by the Department of Energy (DOE), the National Science Foundation, the Jeffress Memorial Trust, the Italian Istituto Nazionale di Fisica Nucleare, the French Institut National de Physique Nucléaire et de Physique des Particules/CNRS and the French Commissariat à l’Energie Atomique. The Southeastern Universities Research Association operates the Thomas Jefferson National Accelerator Facility for the DOE under contract No. DE-AC05-84ER40150.
warning/0506/gr-qc0506128.html	ar5iv	text	# Nonsingular Black Holes and Degrees of Freedom in Quantum Gravity ## Singularities. The main problem caused by a singularity is the fact that it presents a boundary to physical evolution. In order to see whether it persists in quantum gravity, then, the following steps have to be performed. This has to be done in a manner which is independent of coordinate or other gauge choices, and only potential simplifications resulting from the symmetry reduction should be used. One first has to locate classical singularities on the phase space of physical fields, the spatial metric $`q_{ab}`$ and extrinsic curvature related to $`\dot{q}_{ab}`$. Conditions to specify the singular part of phase space must be chosen such that any solution to the theory, which is a trajectory on phase space, intersects this singular part exactly when it develops a singularity. The solution space is in general quite complicated to study, but one can select a variable $`T`$ on phase space which is transversal to the singular part, a local internal time rather than coordinate time, i.e. which fulfills $`T=0`$ in a neighborhood around zero exactly at the singular part. Finally, one needs to write down the quantum evolution of geometry in the local internal time and check whether or not it stops at $`T=0`$. If one can find a $`T`$ such that the quantum evolution does not stop anywhere, the quantum system is non-singular. This is the analog of the classical notion of space-time completeness. We illustrate this scheme with isotropic cosmology where the phase space is 2-dimensional with the scale factor $`a`$ (the spatial radius of the universe) and its time derivative. Singularities occur only if the scale factor vanishes such that $`a=0`$ specifies the singular part. An obvious local internal time (which in this case is global) is given by $`T=a`$, or with a slight modification the triad variable $`p`$ with $`\|p\|=a^2`$ and $`sgnp`$ being the orientation of space. Using this variable makes no difference classically, but is important in quantum geometry where triads are basic variables. At the quantum level one can then first note that operators for $`p^1`$ are finite InvScale , indicating already that curvatures and energy densities do not diverge, and most importantly that the quantum evolution is given by a difference equation for the wave function in $`p`$ which does not stop at $`p=0`$ Sing . Thus, there is no singularity in isotropic loop quantum cosmology. ## Spherical symmetry. The case of interest here is spherical symmetry, where the kinematical phase space on which we have to locate singularities is infinite-dimensional and spanned by the metric components in $$\mathrm{d}s^2=q_x(x)\mathrm{d}x^2+q_\mathrm{\Omega }(x)\mathrm{d}\mathrm{\Omega }^2$$ (1) (in polar coordinates) and their time derivatives. As an example we can look at the Schwarzschild solution for a black hole of mass $`M`$, $`q_x=(12M/x)^1`$, $`q_\mathrm{\Omega }=x^2`$. The singularity is reached for $`x=0`$, at which point both metric coefficients are zero. The question then arises which one, or both, of them must be zero as a condition for a singularity. It turns out that $`q_\mathrm{\Omega }`$ is zero only at the singularity, while $`q_x`$ can also become zero elsewhere, i.e. at the horizon $`x=2M`$, when one chooses a different gauge (e.g. with homogeneous coordinates in the interior). This point illustrates why gauge independence is essential in answering the singularity problem: even the very first step, finding where singularities would develop, depends on it. In fact, in this case we can choose our coordinates $`x`$ and $`t`$ at will, which affects the form of $`q_x`$ and points where it can be zero. In spherical symmetry, however, the fact that $`q_\mathrm{\Omega }=0`$ at the singularity is unaffected (even though, of course, $`q_\mathrm{\Omega }`$ can change as a function of $`x`$ when we change coordinates). We can now consider a spatial slice which locally, around a point $`x_0`$, approaches the classical singularity such that $`q_\mathrm{\Omega }(x_0)0`$. The above discussion shows that $`T=q_\mathrm{\Omega }(x_0)`$ is a good local internal time, which completes setting up the problem from the classical side. It now remains to formulate the quantum evolution in local internal time and check if it stops at $`T=0`$. ## Quantum geometry. Again we first transform to triad variables $`\|E^x\|=q_\mathrm{\Omega }`$ and $`E^\phi =\sqrt{q_xq_\mathrm{\Omega }}`$ which become basic operators in quantum geometry. The (local) orientation of space around a point $`x_0`$ is now given by $`sgnE^x(x_0)`$ where $`E^x`$ unlike $`E^\phi `$ can take both signs. Moreover, the discussion in metric variables shows that $`T=E^x(x_0)`$ is our local internal time such that the situation, so far, is analogous to that in the isotropic case: triad variables lead to a local internal time which takes values at both sides of the classical singularity, $`T=0`$ defining a manifold in superspace rather than at the boundary. It is important to note that the introduction of triad variables was seen as a necessary step toward a background independent quantization. Now it turns out that this also changes the singularity structure on phase space in a way which was important for removing cosmological singularities. Nevertheless, even though the singularity is now located in superspace, the classical evolution still stops there and is not able to connect from positive to negative $`T`$. This still has to be checked by the quantum evolution, the most crucial point. Quantum evolution follows from the Hamiltonian constraint operator acting on states in the form of a lattice model with basis SphSymm $`\|\stackrel{}{k},\stackrel{}{\mu }:=`$ where the integer labels $`k_e`$ on edges are eigenvalues of the operator $`\widehat{E}^x`$ and the positive real labels $`\mu (v)`$ at vertices those of $`\widehat{E}^\phi `$. Positions of vertices do not refer to a background space, and the lattice model represents the continuum theory. The constraint then acts by SphSymmHam $`\widehat{H}[N]`$ $`=_vN(v)(\widehat{C}_0(k)`$ $`+\widehat{C}_{R+}(k)`$ $`+\widehat{C}_R(k)`$ $`+\widehat{C}_{L+}(k)`$ $`+\widehat{C}_L(k)`$ $`+\mathrm{})`$ summing over all vertices of the lattice, the dots indicating further terms such as a matter Hamiltonian whose detailed form is not important here. The known coefficients $`\widehat{C}_I(k)=C_I(k)\widehat{C}_I`$ consist of functions $`C_I(k)`$ of the edge labels and operators $`\widehat{C}_I`$ acting only on the dependence on vertex labels $`\mu `$. A general state is now a superposition $`\|\psi =_{\stackrel{}{k},\stackrel{}{\mu }}\psi (\stackrel{}{k},\stackrel{}{\mu })\|\stackrel{}{k},\stackrel{}{\mu }`$ whose coefficients $`\psi (\stackrel{}{k},\stackrel{}{\mu })`$ define the state in the triad representation. The constraint $`\widehat{H}[N]\|\psi =0`$ has to hold true for all functions $`N`$ with independent values $`N(v)`$, giving one equation for each vertex which in the triad representation takes the form $`\widehat{C}_0(k)\psi (k_{},k_+)+\widehat{C}_{R+}(k)\psi (k_{},k_+2)`$ $`+\widehat{C}_R(k)\psi (k_{},k_++2)+\widehat{C}_{L+}(k)\psi (k_{}2,k_+)`$ $`+\widehat{C}_L(k)\psi (k_{}+2,k_+)+\mathrm{}=0`$ of a difference equation, where we have suppressed the vertex labels on which the $`\widehat{C}_I`$ act and unchanged $`k`$. We now solve this set of equations with initial and boundary values for the wave function. To define a solution scheme we proceed iteratively from vertex to vertex, starting at one side $``$ of the lattice. We assume that the boundary values for all $`\mu _{}`$ and $`k_+()=:k_{}`$ of the wave function as well as values for large positive $`k_e=k_0`$ and $`k_01`$ at all edges $`e`$ are given, which means that we have specified the initial situation, e.g. by a semiclassical state specifying the initial slice far away from the singularity. The equation can then be solved for $`\widehat{C}_{R+}\psi (k_{},k_+2)`$ in terms of values of the wave function specified by the initial conditions. This brings us one step further because we now have information about the wave function at $`k_+2`$ for a smaller edge label (our local internal time) evolving toward the classical singularity. Next, we have to know how to find $`\psi `$ from its image under $`\widehat{C}_{R+}`$. This can be done by specifying conditions for the wave function at small $`\mu `$ (which is not in the singular part of minisuperspace but represents an ordinary boundary) and happens in exactly the same way as in homogeneous models HomCosmo . Before continuing, we notice that this indicates the presence of aspects of the BKL picture in quantum gravity. However, we still have to try to evolve through the classical singularity, i.e. $`k_e=0`$, which will be the main test. One crucial difference to cosmological models is that the coefficients $`\widehat{C}_I(k)`$ are not only functions of the local internal time, $`k_+`$, studied in the iteration but also of neighboring labels such as $`k_{}`$ which do not take part in this difference equation but the dependence on which has been determined in iteration steps for previous vertices. This is clearly a new feature coming from the inhomogeneous context, and it has a bearing on the singularity issue. Singularities are removed if the difference equation determines the wave function everywhere on minisuperspace once initial and boundary conditions have been chosen away from classical singularities. The simplest realization is by a difference equation with non-zero coefficients everywhere. However, this is not automatically the case with an equation coming from a general construction of the Hamiltonian constraint, and so has to be checked explicitly. Here, it turns out SphSymmHam that a symmetric constraint indeed leads to non-zero functions $`C_I(k)`$ which then will not pose a problem to the evolution. All values of the wave function, at positive as well as negative $`k`$, are determined uniquely by the difference equations and chosen initial and boundary values. The evolution thus continues through the classical singularity at zero $`k`$: there is no quantum singularity. Other quantization choices can lead to quantum singularities, providing selection criteria to formulate the quantum theory with implications also for the full framework. ## Consequences. We have shown that the same mechanism as in homogeneous models contributes to the removal of spherically symmetric classical singularities. Key features are that densitized triads as basic variables in quantum geometry provide us with a local internal time taking values at two sides of the classical singularity, combined with a quantum evolution that connects both sides. No new ingredients are necessary for inhomogeneous singularities, only an application of the general scheme to the new and more complicated situation. As in cosmological models the argument applies only to space-like singularities such as the Schwarzschild one. The reason is that we evolve a spatial slice toward the classical singularity and test whether it will stop. A time-like or null singularity would require a different mechanism which is not known at present. Thus, cases like negative mass solutions seem to remain singular, which is a welcome property helping to rule out unwanted solutions leading to instability SingValue . This scenario and its form of difference equations does not only apply to vacuum black holes but also to spherically symmetric matter systems. In such a case, there would be new labels for matter fields, and a contribution to the constraint from the matter Hamiltonian. As this does not change the structure of the difference equation, the same conclusions apply. Moreover, models for Einstein–Rosen waves have a similar structure just with a new vertex label. Also in this case, with or without matter fields, the analysis goes through such that the absence of singularities can be demonstrated even in situations with local gravitational degrees of freedom. There are differences between homogeneous models and these inhomogeneous cases, and the inhomogeneous analysis is much more non-trivial. In homogeneous models there are several ambiguities in the constraint operator, and several choices lead to non-singular evolution. In more complicated situations such as those studied here, not all options remain available. In particular, we had to use a symmetric ordering of the constraint in order to have non-vanishing coefficients of the difference equation. In homogeneous models one can also work with a version whose coefficients vanish right at the singularity. The evolution then still continues since the value at the classical singularity simply decouples and does not play a role for the evolution. Instead, one can use the behavior to find dynamical initial conditions DynIn . This is also possible here for evolution in local internal time, but then the decoupled value at $`k_{}=0`$ is not determined and in general needed for the wave function at other values of $`k_+`$. The inhomogeneous evolution would thus break down, and this choice of constraint is ruled out. There is a difference in the constraint operator we used compared to a common expression in the full theory QSDI . This issue is visible only in inhomogeneous models, and consists in whether or not the constraint creates new edges and vertices, or just changes labels of existing ones. We chose the second possibility, which has already been considered as a modification in the full theory S:ClassLim . There, it can better explain the presence of correlations at an intuitive level, but makes checking anomaly-freedom more complicated. The main problem of an anomalous quantization would be that too many states could be removed when imposing the constraints, leaving not enough physical solutions. This issue can be checked here with the constraint we used. If there is no matter field present we expect just one physical degree of freedom, the Schwarzschild mass $`M`$. In our solution scheme we started with a boundary state $`\psi _{}`$ corresponding to this degree of freedom, and with this state being free it is already clear that we do not lose too many states. It is even possible to check whether or not the number of independent physical solutions is correct, i.e. not too large either. In the iteration we solve one difference equation for $`\psi `$ at each vertex, such that any freedom here would provide new quantum degrees of freedom. Since the difference equation for $`\psi `$ has the same form as that in homogeneous loop quantum cosmology, the number of quantum degrees of freedom is formally related to the initial value problem of quantum cosmology. A possible physical meaning is to be checked in explicit examples. In the isotropic case there are indeed dynamical initial conditions following from the dynamical law DynIn ; Essay which, if realized in our context, would imply that solutions for $`\psi `$ are unique and the mass is the only quantum degree of freedom. However, these conditions rely on the fact that leading coefficients of the difference equation can vanish, which we have ruled out for inhomogeneous models. Moreover, the uniqueness of a quantum cosmological wave function depends on the pre-classicality condition of DynIn . Other mechanisms to select unique cosmological solutions are thus needed, such as from observables or the physical inner product IsoSpinFoam . This issue is quite complicated for difference equations in particular in anisotropic models GenFuncBI . Nonetheless, a simple counting of free variables supports the connection to initial conditions: The vacuum spherically symmetric case has difference equations in three independent variables, an edge label $`k`$ and two neighboring vertex labels $`\mu `$. Homogeneous loop quantum cosmology gives rise to an equation of similar structure and also three variables, so if we assume that there is a mechanism for a unique solution it will also apply to black holes of a given mass. Adding matter fields (or more gravitational freedom as in Einstein–Rosen) increases the number of independent variables to five in inhomogeneous models (two new vertex labels) as opposed to four in homogeneous matter models. The type of difference equations thus agrees in homogeneous and inhomogeneous models in vacuum, but not when local degrees of freedom are present. Thus, the structure of the Hamiltonian constraint equation from loop quantum gravity can potentially provide explanations for issues as diverse as the singularity problem in cosmology and black hole physics, initial conditions in quantum cosmology, the semiclassical limit and issue of quantum degrees of freedom. We emphasize that many of these connections still have to be checked in generality. Still, such connections between seemingly unrelated issues in quantum gravity can be seen as support for the internal consistency of the whole theory and, hopefully, provide guidance for future developments. We can finally come back to the approach to a classical singularity and the BKL picture. Our results here do not rely on an extension of the BKL picture to the quantum situation. First of all, the situation is conceptually different because evolution is now studied for a wave function in local internal time $`T`$, rather than the spatial metric in coordinate time. Nevertheless, at first sight a similar picture arises here from the quantum equation: as used in the previous arguments, the equations can be reduced to ordinary difference equations in $`T`$, where neighboring edges just contribute via an inhomogeneity of the difference equation. The inhomogeneous situation, however, does play an important role right at the classical singularity where some versions which would be allowed in homogeneous models are ruled out. Given that the techniques necessary for the quantum theory are similar to lattice models, it is easy to implement them in numerical quantum gravity. This opens the door to numerical investigations of many problems that are still actively pursued in classical gravity NumSing , such as the approach to classical singularities and the issue of gravitational collapse and naked singularities. This requires studying horizons in addition to classical singularities, which can also be done at the quantum level Horizon . As we have seen, there are many non-trivial quantum effects which play together in just the right way to ensure the absence of singularities, which has prospects for other effects in the physics of black holes BHPara .
warning/0506/math0506291.html	ar5iv	text	# Galois Corings and a Jacobson-Bourbaki type Correspondence. ## Introduction One of the key pieces in the Galois theory of fields and more generally of division rings is the Jacobson-Bourbaki Theorem, see \[11, Chapter 7, Sections 2, 3\] and \[10, Section 8.2\]. Let $`E`$ be a division ring with prime field $`k`$. Consider the injective ring homomorphism $`r:E\mathrm{End}({}_{k}{}^{}E),er_e`$ where $`r_e(e^{})=e^{}e`$ for all $`e^{}E`$ (the multiplication in $`\mathrm{End}({}_{k}{}^{}E)`$ is the opposite of the composition). The Jacobson-Bourbaki Theorem states that there is a bijective correspondence between the set of division subrings $`D`$ of $`E`$ such that $`{}_{D}{}^{}E`$ is finite dimensional and the set of subrings $`S`$ of $`\mathrm{End}({}_{k}{}^{}E)`$ such that $`Im(r)S`$ and $`S_E`$ is finite dimensional. The ring $`\mathrm{End}({}_{k}{}^{}E)`$ is indeed an $`E`$-ring and the condition $`Im(r)S`$ can be rephrased as $`S`$ being an $`E`$-subring of $`\mathrm{End}({}_{k}{}^{}E)`$. This correspondence is hidden behind the veil of the Galois connection in a Galois extension of fields or more generally of division rings, see and . Using the dual structure of $`E`$-ring, the structure of $`E`$-coring, in Sweedler gave a dual result to the Jacobson-Bourbaki Theorem. The advantage of using this dual structure is that the finiteness conditions needed in the Jacobson-Bourbaki Theorem can be dropped. The finiteness conditions come implicit in the structure of coring through the fact that an element is mapped into a finite sum of elements via the comultiplication map. Sweedler’s result asserts that there is a bijective correspondence between division subrings of $`E`$ and quotient corings of the $`E`$-coring $`E_kE`$. The Jacobson-Bourbaki Theorem can be obtained from Sweedler’s result by duality and this process makes clear why the finiteness conditions are needed. The goal of this paper is to extend Sweedler’s result from division rings to simple artinian rings replacing the Sweedler coring $`E_kE`$ by a more general type of coring, the comatrix coring introduced in to describe the structure of cosemisimple corings. Let $`\mathrm{\Sigma }`$ be a finitely generated and projective right module over a ring $`A`$ and let $`B`$ be a simple artinian subring of $`End(\mathrm{\Sigma }_A)`$. Let $``$ denotes the comatrix $`A`$-coring $`\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$ constructed on the bimodule $`\mathrm{\Sigma }`$ and with coefficients $`B`$, see (1). Our main theorem (Theorem 2.4) states that there is a bijective correspondence between the set of all simple artinian subrings $`BCEnd(\mathrm{\Sigma }_A)`$ and the set of all coideals $`J`$ of $``$ such that the quotient coring $`/J`$ is simple cosemisimple. If in addition $`\mathrm{\Sigma }_A`$ is simple, then any quotient coring of $``$ is simple cosemisimple, thus obtaining a bijective correspondence between intermediate division subrings of $`B\mathrm{End}(\mathrm{\Sigma }_A)`$ and coideals of $``$ (see Remark 2.5). Sweedler’s result, together with some additional information on conjugated subextensions, is then obtained as a consequence by taking $`A`$ a division ring and $`\mathrm{\Sigma }=A`$ (Corollary 2.6). An example illustrating the bijective correspondence is worked out. This closes Section 2, which thus contains our main results. Section 1 recalls the most fundamental results on comatrix corings, Galois corings, and cosemisimple corings needed in the sequel. We also include a homological characterization of Galois corings (Theorem 1.5), which gives as a consequence that if the canonical map is surjective for a quasi-projective comodule with a generating condition, then the coring is Galois (Corollary 1.6). This corollary is used in the proof of the main result of Section 3, namely, Theorem 3.2, which states a Jacobson-Bourbaki correspondence for simple artinian subextensions of a ring extension. This correspondence is dual to the stated in Section 2. We complete the paper with an Appendix that contains a complete classification of the simple cosemisimple $`/`$–corings. We next fix notation and present some basic definitions. In the sequel $`A,B`$ denote associative and unitary algebras over a commutative ring $`K`$. By $`_A`$ we denote the tensor product over $`A`$. The category of right $`A`$-modules is denoted by $`_A`$. Bimodules are assumed to be centralized by $`K`$. An $`A`$-coring (or $`A/K`$–coring, when $`K`$ is not obvious from the context) is a triple $`(,\mathrm{\Delta },ϵ)`$ where $``$ is an $`A`$-bimodule and $`\mathrm{\Delta }:_A`$ (comultiplication) and $`ϵ:A`$ (counit) are $`A`$-bimodule maps such that $`(id_{}_A\mathrm{\Delta })\mathrm{\Delta }=(\mathrm{\Delta }_Aid_{})\mathrm{\Delta }`$ and $`(ϵ_Aid_{})\mathrm{\Delta }=(id_{}_Aϵ)\mathrm{\Delta }=id_{}`$. In a categorical language, a coring is just a coalgebra in the monoidal category of $`A`$-bimodules with the tensor product $`_A`$ as a product. For $`c`$ we will write $`\mathrm{\Delta }(c)=c_{(1)}_Ac_{(2)}`$. The left dual $`{}_{}{}^{}=Hom({}_{A}{}^{},{}_{A}{}^{}A)`$ of the coring $``$ is an $`A^{opp}`$-ring ($`A^{opp}`$ denotes the opposite ring of $`A`$) with the product $`(fg)(c)=f(c_{(1)}g(c_{(2)}))`$ for all $`f,g{}_{}{}^{}`$ and $`c`$. Similarly, the right dual $`^{}`$ of $``$ is an $`A^{opp}`$-ring with the product $`(fg)(c)=g(f(c_{(1)})c_{(2)})`$ for $`f,g^{}`$ and $`c`$. A right $``$-comodule is a right $`A`$-module together with an $`A`$-module map $`\rho _M:MM_A`$ such that $`(id_M_A\mathrm{\Delta })\rho _M=(\rho _Mid_{})\rho _M`$ and $`(id_Mϵ)\rho _M=id_M`$. A $``$-comodule map between two right $``$-comodules $`M`$ and $`N`$ is an $`A`$-module map $`f:MN`$ such that $`(f_Aid_{})\rho _M=\rho _Nf`$. By $`Hom_{}(M,N)`$ we will denote the $`K`$-module of all $``$-comodule maps between $`M`$ and $`N`$. The category whose objects are right $``$-comodules and whose morphisms are $``$-comodule maps is denoted by $`^{}`$. It is an additive $`K`$–linear category and if $`{}_{A}{}^{}`$ is flat, then it is a Grothendieck category. The product of every endomorphism ring of an object in an additive category is by default the composition. We adopt, however, the following convention in the case of modules: the product of the endomorphism ring $`\mathrm{End}(M_A)`$ of a right $`A`$–module $`M`$ is the composition, although by $`\mathrm{End}({}_{A}{}^{}N)`$ we will denote the opposite ring of the endomorphism ring of a left $`A`$–module $`N`$, being then its product the opposite of the composition. ## 1 Comatrix corings, Galois comodules and cosemisimple corings Let $`\mathrm{\Sigma }`$ be a $`BA`$–bimodule, and assume that $`\mathrm{\Sigma }_A`$ is finitely generated and projective with a finite dual basis $`\{(e_i^{},e_i)\}\mathrm{\Sigma }^{}\times \mathrm{\Sigma }`$. We can consider a coring structure \[6, Proposition 2.1\] over the $`A`$–bimodule $`\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$ with comultiplication and counit defined respectively by $$\mathrm{\Delta }(\varphi _Bx)=\underset{i}{}\varphi _Be_i_Ae_i^{}_Bx,ϵ(\varphi _Bx)=\varphi (x).$$ (1) The comultiplication is independent of the choice of the dual basis. This coring will be called the $`A`$-comatrix coring on $`\mathrm{\Sigma }`$ with coefficients in $`B`$. The $`A`$-module $`\mathrm{\Sigma }`$ becomes a right $`\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$–comodule with coaction $$\varrho _\mathrm{\Sigma }(x)=\underset{i}{}e_i_Ae_i^{}_Bx.$$ Assume $`\mathrm{\Sigma }`$ to be the underlying $`A`$–module of a right comodule over some $`A`$–coring $``$, with structure map $`\rho _\mathrm{\Sigma }:\mathrm{\Sigma }\mathrm{\Sigma }_A`$. In this case, with $`T=\mathrm{End}(\mathrm{\Sigma }_{})`$, we have from \[6, Proposition 2.7\] that the map $`\mathrm{𝖼𝖺𝗇}:\mathrm{\Sigma }^{}_T\mathrm{\Sigma }`$ defined by $$\mathrm{𝖼𝖺𝗇}(\varphi _Tx)=\varphi (x_0)x_1(\rho _\mathrm{\Sigma }(x)=x_0_Ax_1)$$ is a homomorphism of $`A`$–corings. This canonical map allowed to extend \[6, Definition 3.4\] the notion of a Galois coring without assuming the existence of group-like elements. When the role of the comodule $`\mathrm{\Sigma }`$ is stressed, the terminology of Galois comodules introduced in is more convenient. Probably, the best solution here is to mention both the coring and the comodule. ###### Definition 1.1. The pair $`(,\mathrm{\Sigma })`$ is said to be Galois if $`\mathrm{𝖼𝖺𝗇}`$ is an isomorphism. In such a case, we say that $``$ is a Galois coring and $`\mathrm{\Sigma }`$ is termed a Galois comodule. The extension $`T\mathrm{End}(\mathrm{\Sigma }_{})`$ is called a $`(,\mathrm{\Sigma })`$-Galois extension. The notion of a noncommutative $`G`$–Galois extension, may be recovered from this definition, see \[6, Example 2.9\]. In this case, the corresponding Galois coring has a group-like element. We give an example of a noncommutative Galois extension for a coring without group-like elements. ###### Example 1.2. Let $``$ and $``$ denote the complex number field and the Hamilton’s quaternions algebra respectively. Consider the right $``$-vector space $`𝔗`$ with basis $`\{c,s\}`$. This vector space becomes a $``$-bimodule with left action $`ic=c`$ and $`is=s`$ and a $``$-coring with comultiplication and counity defined by $$\begin{array}{cc}\mathrm{\Delta }(c)=ccss,\hfill & ϵ(c)=1,\hfill \\ \mathrm{\Delta }(s)=cs+sc,\hfill & ϵ(s)=0.\hfill \end{array}$$ Analogously to the coalgebra case this coring is called the trigonometric coring. This coring has no group-like elements. Let $`\mathrm{\Sigma }`$ be a right $``$-vector space with basis $`\{v_1,v_2\}`$. The map $`\rho :\mathrm{\Sigma }\mathrm{\Sigma }𝔗`$ defined by $$v_1v_1c+v_2s,v_2v_2cv_1s$$ makes $`\mathrm{\Sigma }`$ into a right $`𝔗`$-comodule. It is not difficult to check that $`(𝔗,\mathrm{\Sigma })`$ is a Galois $``$-coring. The corresponding Galois extension is the well-known embedding of $``$ into $`M_2()`$: $$i\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),j\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right).$$ Our first objective is to enrich \[3, 18.26\] with a new characterization of Galois comodules. We need two previous observations. The first one is that if $``$ is flat as a left $`A`$–module then, using that the forgetful functor $`U:^{}_A`$ is faithful and exact, the following lemma can be proved (see ). ###### Lemma 1.3. If $``$ is flat as a left $`A`$–module, then a right $``$–comodule $`M`$ is finitely generated in the Grothendieck category $`^{}`$ if and only if $`M`$ is finitely generated as a right $`A`$–module. We have a pair of functors (2) where $`_T\mathrm{\Sigma }`$ is left adjoint to $`\mathrm{Hom}_{}(\mathrm{\Sigma },)`$. If $`\chi :\mathrm{Hom}_{}(\mathrm{\Sigma },)_T\mathrm{\Sigma }id_{^{}}`$ is the counit of this adjunction, then the canonical map can be expressed \[6, Lemma 3.1\] as the composite (3) Observe that $`\chi _{}`$ is an isomorphism if and only if $`\mathrm{𝖼𝖺𝗇}`$ is so. Our second observation is the following lemma. ###### Lemma 1.4. Let $`\eta :id__T\mathrm{Hom}_{}(\mathrm{\Sigma },_T\mathrm{\Sigma })`$ be the unit of the adjunction (2). Consider the map $`\mathrm{𝖼𝖺𝗇}_{}:\mathrm{Hom}_{}(\mathrm{\Sigma },\mathrm{\Sigma }^{}_T\mathrm{\Sigma })\mathrm{Hom}_{}(\mathrm{\Sigma },),f\mathrm{𝖼𝖺𝗇}f.`$ Then the following composition is the identity map on $`\mathrm{\Sigma }^{}`$ ###### Proof. If we apply the displayed composite map to $`\varphi \mathrm{\Sigma }^{}`$, then we obtain the map from $`\mathrm{\Sigma }`$ to $`A`$ given by $`xϵ_{}(\varphi (x_0)x_1)`$, which is nothing but $`\varphi `$ since $$ϵ_{}(\varphi (x_0)x_1)=\varphi (x_0)ϵ_{}(x_1)=\varphi (x_0ϵ_{}(x_1))=\varphi (x)$$ We are now in position to state our homological characterization of Galois comodules. ###### Theorem 1.5. Let $`\mathrm{\Sigma }`$ be a right comodule over an $`A`$–coring $``$ and assume that $`\mathrm{\Sigma }`$ is finitely generated and projective as a right $`A`$-module. If $`{}_{A}{}^{}`$ is flat, then $`(,\mathrm{\Sigma })`$ is Galois if and only if there exists an exact sequence $`\mathrm{\Sigma }^{(J)}\mathrm{\Sigma }^{(I)}0`$ in $`^{}`$ such that the sequence is exact. ###### Proof. Consider $`T^{(J)}T^{(I)}\mathrm{\Sigma }^{}0`$ a free presentation of the left $`T`$–module $`\mathrm{\Sigma }^{}`$. By tensorizing with $`{}_{T}{}^{}\mathrm{\Sigma }`$ we obtain an exact sequence $`\mathrm{\Sigma }^{(J)}\mathrm{\Sigma }^{(I)}\mathrm{\Sigma }^{}_T\mathrm{\Sigma }0.`$ We have then the following commutative diagram in $`_T`$ (4) Now, $`\eta _{T^{(J)}}`$ and $`\eta _{T^{(I)}}`$ are isomorphisms because $`\mathrm{\Sigma }`$ is finitely generated in the Grothendieck category $`^{}`$ (see ). By Lemma 1.4, $`\eta _\mathrm{\Sigma }^{}`$ is an isomorphism if and only if $`\mathrm{𝖼𝖺𝗇}_{}`$ is bijective. So if we assume that $`(,\mathrm{\Sigma })`$ is Galois, then $`\eta _\mathrm{\Sigma }^{}`$ is an isomorphism and, from (4), the following sequence is exact Observe that the isomorphism of $`A`$–corings $`\mathrm{𝖼𝖺𝗇}:\mathrm{\Sigma }^{}_T\mathrm{\Sigma }`$ is also an isomorphism of right $``$–comodules. This finishes the proof of the necessary condition. For the sufficiency, consider the commutative diagram in $`^{}`$ with exact rows In this diagram, $`\chi _{\mathrm{\Sigma }^{(J)}}`$ and $`\chi _{\mathrm{\Sigma }^{(I)}}`$ are isomorphisms because $`\mathrm{\Sigma }`$ is finitely generated in $`^{}`$. We have then that $`\chi _{}`$ is an isomorphism and, thus, $`(,\mathrm{\Sigma })`$ is Galois. ∎ T. Brzeziński has shown in that a simple comodule with $`\mathrm{𝖼𝖺𝗇}`$ surjective is Galois. As a consequence of Theorem 1.5 we derive a generalization of Brzeziński’s result. Following the definition given for modules in , we say the comodule $`\mathrm{\Sigma }`$ is quasi-projective if for every exact sequence $`\mathrm{\Sigma }N0`$ in $`^{}`$ then the sequence of abelian groups $`\mathrm{Hom}_{}(\mathrm{\Sigma },\mathrm{\Sigma })\mathrm{Hom}_{}(\mathrm{\Sigma },N)0`$ is exact. Since, by Lemma 1.3, $`\mathrm{\Sigma }`$ is finitely generated in $`^{}`$, a straightforward adaptation of \[1, Proposition 16.2.(2)\] to Grothendieck categories gives that $`\mathrm{Hom}_{}(\mathrm{\Sigma },)`$ will already preserve exact sequences of the form $`\mathrm{\Sigma }^{(I)}N0`$. The following corollary is then easily deduced from Theorem 1.5, making use once more of the “AB5” condition. ###### Corollary 1.6. Assume that $`\mathrm{\Sigma }`$ is quasi-projective in $`^{}`$ and generates every subcomodule of any finite direct sum of copies of $`\mathrm{\Sigma }`$ (e.g., $`\mathrm{\Sigma }`$ is a semisimple comodule). If $`\mathrm{𝖼𝖺𝗇}`$ is surjective, then $`(,\mathrm{\Sigma })`$ is Galois. Now, let us recall from the structure of cosemisimple corings, which is tightly related to the coring version of the Generalized Descent Theorem formulated in \[6, Theorem 3.10\]. A coring is said to be cosemisimple if it satisfies one of the equivalent conditions of the following theorem. ###### Theorem 1.7. \[7, Theorem 3.1\] The following assertions for an $`A`$-coring $``$ are equivalent: 1. Every left $``$-comodule is semisimple and $`{}_{}{}^{}`$ is abelian. 2. Every right $``$-comodule is semisimple and $`^{}`$ is abelian. 3. $`{}_{}{}^{}`$ is semisimple and $`_A`$ is flat. 4. $`_{}`$ is semisimple and $`{}_{A}{}^{}`$ is flat. A coring is called simple if it has no non trivial subbicomodules. It was proved in \[7, Theorem 3.7\] that any cosemisimple coring decomposes in a unique way as a direct sum of simple cosemisimple corings. An example of simple cosemisimple coring is the comatrix $`A`$-coring $`\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$, where $`\mathrm{\Sigma }_A`$ is finitely generated and projective and $`BEnd(\mathrm{\Sigma }_A)`$ is simple artinian. The following result shows that indeed all simple cosemisimple corings can be obtained in this way. ###### Proposition 1.8. \[6, Proposition 4.2\] Let $``$ be a simple cosemisimple $`A`$-coring and $`\mathrm{\Sigma }_{}`$ a finitely generated right $``$-comodule. Then $`T=End(\mathrm{\Sigma }_{})`$ is simple artinian, $`\mathrm{\Sigma }_A`$ is finitely generated projective and the canonical map $`\mathrm{𝖼𝖺𝗇}:\mathrm{\Sigma }^{}_T\mathrm{\Sigma }`$ is an isomorphism. A more precise description of simple cosemisimple corings is given by the following structure theorem. It may be viewed as a generalization of the Artin-Wedderburn Theorem. ###### Theorem 1.9. \[6, Theorem 4.3\] An $`A`$-coring $``$ is simple cosemisimple if and only if there is a finitely generated projective right $`A`$-module $`\mathrm{\Sigma }`$ and a division subring $`DEnd(\mathrm{\Sigma }_A)`$ such that $`\mathrm{\Sigma }^{}_D\mathrm{\Sigma }`$ as $`A`$-corings. In such a case, if $`\mathrm{\Gamma }`$ is another finitely generated and projective right $`A`$-module and $`EEnd(\mathrm{\Gamma }_A)`$ is a division subring, then $`\mathrm{\Gamma }^{}_E\mathrm{\Gamma }`$ if and only if there is an isomorphism of right $`A`$-modules $`g:\mathrm{\Sigma }\mathrm{\Gamma }`$ such that $`gDg^1=E.`$ In view of this structure theorem, for a field extension $`A/k`$, the classification of simple cosemisimple $`A`$-corings centralized by $`k`$ is reduced to the classification of finite dimensional division algebras over $`k`$ and to the study of how these division algebras embed in matrix algebras over $`A`$. The first problem leads to the Brauer group theory of a field and the second one can be treated with the help of the Skolem-Noether Theorem. The complete classification for the field extension $`/`$ is obtained in the Appendix. ## 2 The Galois connection from a coring point of view Let $`\mathrm{\Sigma }_A`$ be a finitely generated and projective right $`A`$-module and denote by $`S=\mathrm{End}(\mathrm{\Sigma }_A)`$ its endomorphism ring. Let $`BS`$ be a subring, and consider the comatrix $`A`$–coring $`=\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$. A typical situation is to consider $`B=k`$, a field, and $`S=M_n(k)`$, the ring of square matrices of order $`n`$ over $`k`$. Let $`\mathrm{𝖲𝗎𝖻𝖾𝗑𝗍}(S/B)`$ denote the set of all ring subextensions $`BCS`$, and $`\mathrm{𝖢𝗈𝗂𝖽𝖾𝖺𝗅𝗌}()`$ be the set of all coideals of $``$. Consider the maps (5) defined as follows. For each subextension $`C\mathrm{𝖲𝗎𝖻𝖾𝗑𝗍}(S/B)`$ we have a canonical homomorphism of $`A`$–corings $`=\mathrm{\Sigma }^{}_B\mathrm{\Sigma }\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$, whose kernel $`𝒥(C)`$ is a coideal of $``$. Conversely, given a coideal $`J`$ of $``$, then, by \[6, Proposition 2.5\], $`B\mathrm{End}(\mathrm{\Sigma }_{})\mathrm{End}(\mathrm{\Sigma }_{/J})S`$, so that $`(J)=\mathrm{End}(\mathrm{\Sigma }_{/J})`$ is a subextension of $`BS`$. Both $`𝒥`$ and $``$ are inclusion preserving maps. The following proposition, which generalizes \[13, Proposition 6.1\], collects some more of their relevant general properties. ###### Proposition 2.1. The maps defined in (5) enjoy the following properties. 1. $`𝒥(C)C`$ for every $`C\mathrm{𝖲𝗎𝖻𝖾𝗑𝗍}(S/B)`$. 2. $`𝒥(J)J`$ for every $`J\mathrm{𝖢𝗈𝗂𝖽𝖾𝖺𝗅𝗌}()`$. 3. $`𝒥(C)=C`$ if and only if $`\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})=C`$. 4. $`𝒥(J)=J`$ if and only if $`(/J,\mathrm{\Sigma })`$ is Galois. 5. The maps $`𝒥`$ and $``$ establish a bijection between the set consisting of ring subextensions $`BCS`$ such that $`\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})=C`$ and the set of coideals $`J`$ of $``$ such that $`(/J,\mathrm{\Sigma })`$ is Galois. ###### Proof. A direct computation gives that $`𝒥(C)=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})`$ for every $`C\mathrm{𝖲𝗎𝖻𝖾𝗑𝗍}(S/B)`$. Thus, (1) follows from \[6, Proposition 2.5\]. This gives also (3). Now, for a given coideal $`J\mathrm{𝖢𝗈𝗂𝖽𝖾𝖺𝗅𝗌}()`$, we have the commutative diagram with exact rows (6) which implies (2) and (4). Finally, let us prove (5): for any ring subextension $`C\mathrm{𝖲𝗎𝖻𝖾𝗑𝗍}(S/B)`$, the $`A`$-coring $`\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$ is Galois by \[6, Lemma 3.9\]. Hence $`/𝒥(C)`$ is Galois. Conversely, for any $`J\mathrm{𝖢𝗈𝗂𝖽𝖾𝖺𝗅𝗌}()`$ we have $`\mathrm{End}(\mathrm{\Sigma }_{})\mathrm{End}(\mathrm{\Sigma }_{/J})`$. If $`/J`$ is Galois, then the canonical map $`\mathrm{𝖼𝖺𝗇}_{/J}`$ is an isomorphism and, therefore, $`\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_{/J})}\mathrm{\Sigma }})=\mathrm{End}(\mathrm{\Sigma }_{/J})`$. Statement (5) follows now from (3) and (4). ∎ ###### Remark 2.2. If $`{}_{A}{}^{}(\mathrm{\Sigma }^{}_C\mathrm{\Sigma })`$ is locally projective (see e.g. \[3, 42.10\] for this notion), then $`C=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})`$ if and only if $`{}_{C}{}^{}\mathrm{\Sigma }`$ is faithfully balanced, i.e., $`C`$ is isomorphic to the biendomophisms ring of $`{}_{C}{}^{}\mathrm{\Sigma }`$ under the natural map. This is because for $`{}_{A}{}^{}(\mathrm{\Sigma }^{}_C\mathrm{\Sigma })`$ locally projective, $`\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})=\mathrm{End}(_{{}_{}{}^{}(\mathrm{\Sigma }^{}_C\mathrm{\Sigma })}\mathrm{\Sigma })`$, see \[3, 19.3\]. Then, by \[6, Proposition 2.1\], $`{}_{}{}^{}(\mathrm{\Sigma }^{}_C\mathrm{\Sigma })\mathrm{End}({}_{C}{}^{}\mathrm{\Sigma })^{op}`$ canonically and, therefore, $$\begin{array}{cc}\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})\hfill & =\mathrm{End}(_{{}_{}{}^{}(\mathrm{\Sigma }^{}_C\mathrm{\Sigma })}\mathrm{\Sigma })=\mathrm{End}({}_{\mathrm{End}({}_{C}{}^{}\mathrm{\Sigma })^{op}}{}^{}\mathrm{\Sigma })=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{End}({}_{C}{}^{}\mathrm{\Sigma })}).\hfill \end{array}$$ ###### Remark 2.3. We have from \[6, Proposition 2.5\] that $$\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})=\{f\mathrm{End}(\mathrm{\Sigma }_A)\|f_Cx=1_Cf(x),\text{ for every }x\mathrm{\Sigma }\}:=\overline{C}.$$ Since $`\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$ is Galois, we get that $`\overline{\overline{C}}=\overline{C}`$. Proposition 2.1 gives a bijective correspondence between coideals $`J`$ of $``$ such that $`/J`$ is Galois and subextensions $`C`$ such that $`C=\overline{C}`$. We are now ready to state a generalization of Sweedler’s predual to Jacobson-Bourbaki Theorem \[14, Theorem 2.1\]. We will say that $`C,C^{}\mathrm{𝖲𝗎𝖻𝖾𝗑𝗍}(S/B)`$ are conjugated in $`S`$ if there is an unit $`gS`$ such that $`C^{}=gCg^1`$. ###### Theorem 2.4. Let $`\mathrm{\Sigma }`$ be a finitely generated projective right $`A`$-module and let $`S=\mathrm{End}(\mathrm{\Sigma }_A)`$. Let $`B`$ be a subring of $`S`$ and consider the comatrix $`A`$-coring $`=\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$. Denote by $`𝒮`$ the set of all simple artinian subrings $`BCS`$ and let $`𝒯`$ denote the set of all coideals $`J`$ of $``$ such that $`/J`$ is simple cosemisimple. Then the maps $$\begin{array}{c}():𝒯𝒮,J\mathrm{End}(\mathrm{\Sigma }_{/J})\hfill \\ 𝒥():𝒮𝒯,CKer(\mathrm{\Sigma }^{}_C\mathrm{\Sigma })\hfill \end{array}$$ are inverse to each other. If, in addition, $`A`$ is a (noncommutative) local ring, then two intermediate simple artinian subrings $`C`$ and $`C^{}`$ are conjugated in $`S`$ if and only if $`/𝒥(C)`$ and $`/𝒥(C^{})`$ are isomorphic as $`A`$–corings. ###### Proof. For $`C𝒮`$, the $`A`$-coring $`\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$ is a simple cosemisimple $`A`$-coring in virtue of \[6, Proposition 4.2\]. Hence $`/𝒥(C)\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$ is simple cosemisimple and so $`𝒥(C)𝒯`$. From \[6, Proposition 2.5\], $`C\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }}).`$ Since $`{}_{C}{}^{}\mathrm{\Sigma }`$ is faithfully flat, \[6, Theorem 3.10\] yields $`C=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})`$. By Theorem 2.1, $`𝒥(C)=C`$. Assume that $`J𝒯`$, then $`/J`$ is simple cosemisimple. By \[6, Theorem 4.1\], $`/J`$ is flat as a left $`A`$-module which implies, by \[9, Lemma 3.1\], that $`\mathrm{\Sigma }`$ is finitely generated as a right $`/J`$-comodule. Hence $`\mathrm{End}(\mathrm{\Sigma }_{/J})`$ is a simple artinian ring. So $`(J)𝒮`$. Furthermore, $`/J`$ is Galois by \[6, Proposition 4.2\]. By Theorem 2.1, $`𝒥(J)=J`$. For the second assertion, let $`C,C^{}`$ be intermediate simple artinian subrings and let $`gS`$ be invertible such that $`gCg^1=C^{}`$. We check that $`\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$ is isomorphic to $`\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$. The set $`\{e_i^{}g,g(e_i)\}_{i=1}^n`$ is a dual basis for $`\mathrm{\Sigma }`$ and, by \[6, Remark 2.2\] the coring structure on $`\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$ defined by this dual basis is the same as the one defined by $`\{e_i^{},e_i\}_{i=1}^n.`$ Consider the $`A`$-bimodule map, $$\psi :\mathrm{\Sigma }^{}\times \mathrm{\Sigma }\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma },(\phi ,x)\phi g^1_C^{}g(x).$$ Let $`cC`$ and $`c^{}=gcg^1C^{}`$. Then, $$\begin{array}{cc}\psi (\phi c,x)\hfill & =\phi cg^1_C^{}g(x)=\phi g^1c^{}_C^{}g(x)=\phi g^1_C^{}c^{}g(x)\hfill \\ & =\phi g^1_C^{}gc(x)=\psi (\phi ,cx).\hfill \end{array}$$ Hence $`\psi `$ defines a unique $`A`$-bimodule map $`\overline{\psi }:\mathrm{\Sigma }^{}_C\mathrm{\Sigma }\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$. It is routine to verify that $`\psi `$ is an isomorphism of $`A`$-corings. Conversely, let $`\chi :\mathrm{\Sigma }^{}_C\mathrm{\Sigma }\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$ be an isomorphism of $`A`$-corings. Let $`S,S^{}`$ be the unique, up to isomorphism, simple right $`\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$-comodule and $`\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$-comodule, respectively. Then $`\mathrm{\Sigma }S^{(m)}`$ as a $`\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$-comodule and $`\mathrm{\Sigma }S^{(n)}`$ as a $`\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$-comodule for some $`m,n`$. Let $`S^\chi `$ denote $`S`$ viewed as a right $`\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$-comodule via $`\chi `$. Then $`S^\chi S^{}`$ as a $`\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$-comodule. In particular, they are isomorphic as right $`A`$-modules. Since $`A`$ is local, $`S^\chi S^{}A^{(l)}`$ as right $`A`$-modules for a certain $`l`$. Let $`\mathrm{\Sigma }^\chi `$ denote $`\mathrm{\Sigma }`$ when considered as a right $`\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$–comodule. Then $`\mathrm{\Sigma }^\chi (S^\chi )^{(m)}A^{(lm)}`$ as a right $`A`$-module. On the other hand, $`\mathrm{\Sigma }A^{(ln)}`$ as a right $`A`$-module. Since the underlying right $`A`$-module of $`\mathrm{\Sigma }^\chi `$ and $`\mathrm{\Sigma }`$ is the same, $`m=n`$ and hence $`\mathrm{\Sigma }^\chi \mathrm{\Sigma }`$ as a right $`\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }`$-comodule. Denote this isomorphism by $`g`$. Then $`\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }})=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }}^\chi )`$ via $`dg^1dg`$. As $`\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }}^\chi )=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }}),`$ $`C=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C\mathrm{\Sigma }})`$ and $`C^{}=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_C^{}\mathrm{\Sigma }}),`$ the assertion holds. ∎ ###### Remark 2.5. With hypothesis as in Theorem 2.4, if we assume in addition that $`\mathrm{\Sigma }_A`$ is simple, then any quotient coring of $`=\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$ is simple cosemisimple: since $`\mathrm{\Sigma }_A`$ is simple, $`\mathrm{\Sigma }`$ is simple as a right $`/J`$-comodule for any coideal $`J`$ of $``$. It is easy to see that the canonical map $`\mathrm{𝖼𝖺𝗇}:\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_{/J})}\mathrm{\Sigma }/J`$ is surjective. Using that $`\mathrm{\Sigma }_{/J}`$ is simple, by Corollary 1.6, $`\mathrm{𝖼𝖺𝗇}`$ is an isomorphism. As $`\mathrm{End}(\mathrm{\Sigma }_{/J})`$ is a division ring, $`/J`$ is simple cosemisimple. Thus we have a bijective correspondence between intermediate division subrings of $`B\mathrm{End}(\mathrm{\Sigma }_A)`$ and coideals of $``$. Observe that no assumptions are made on $`A`$. As a consequence we can derive Sweedler’s predual to the Jacobson-Bourbaki Theorem, with the additional information concerning conjugated subrings. ###### Corollary 2.6. \[14, Theorem 2.1\] Let $`DE`$ be division rings. Set $`=E_DE`$ and let $`g=1_D1`$ be the distinguished group-like element. For a coideal $`J`$ of $``$ let $`\pi _J:/J`$ denote the canonical projection. Then, the maps $$\begin{array}{c}():𝒯𝒮,J\{eE:e\pi _J(g)=\pi _J(g)e\}\hfill \\ 𝒥():𝒮𝒯,CKer(\mathrm{\Sigma }^{}_C\mathrm{\Sigma })\hfill \end{array}$$ establish a bijective correspondence between the set $`𝒮`$ of intermediate division rings $`DCE`$ and the set $`𝒯`$ of coideals $`J`$ of $``$. Moreover, two intermediate division rings $`C`$ and $`C^{}`$ are conjugated in $`S`$ if and only if $`/𝒥(C)/𝒥(C^{}).`$ ###### Proof. It only remains to prove that $`\mathrm{End}(E_{/J})=\{eE:e\pi _J(g)=\pi _J(g)e\}`$ but this is easily checked. ∎ ###### Remark 2.7. If $`\mathrm{\Sigma }_A`$ is not simple, then the quotient corings of $`\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$ need not in general to be simple cosemisimple. Let $`k`$ be a field, $`\mathrm{\Sigma }=k^{(2)}`$ and $`T=T_2(k)\mathrm{End}({}_{k}{}^{}\mathrm{\Sigma })=M_2(k)`$ the upper triangular matrix algebra. Then $`\mathrm{\Sigma }^{}_T\mathrm{\Sigma }`$ is a non simple cosemisimple quotient coalgebra of $`\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$. This example also serves to show that factor corings of Galois corings are not Galois. The module $`\mathrm{\Sigma }_T`$ is isomorphic to the indecomposable projective $`eT`$, where $`eM_2(k)`$ is the elementary matrix with $`1`$ in the $`(1,1)`$-entry and zero elsewhere. Thus, $`\mathrm{End}(\mathrm{\Sigma }_T)eTek`$, and the canonical map (15) gives here a surjective $`k`$-coalgebra homomorphism $`\mathrm{𝖼𝖺𝗇}:\mathrm{\Sigma }^{}_k\mathrm{\Sigma }T^{}`$ which, obviously, cannot be bijective. Thus, if we take $`=\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$, and $`𝔇=T^{}`$, then the factor coalgebra $`(𝔇,\mathrm{\Sigma })`$ of the Galois coalgebra $`(,\mathrm{\Sigma })`$ is not Galois. On the other hand, observe that $`\mathrm{\Sigma }_T^{}`$ is projective but does not generate all its submodules, which sheds some light on the conditions involved in Theorem 1.5. ###### Example 2.8. We next illustrate the Galois connection established in Theorem 2.4 by a concrete example. Assume that $`k`$ has an $`n`$-th primitive root of unity $`\omega `$. For $`\alpha ,\beta `$ non zero elements in $`k`$ let $`A_\omega (\alpha ,\beta )`$ denote the associative $`k`$-algebra generated by two elements $`x,y`$ subject to the relations $`x^n=\alpha ,y^n=\beta `$ and $`yx=\omega xy`$. Details on the properties of this algebra to be used in the sequel may be consulted in \[12, Chapter 15\]. The algebra $`A_\omega (\alpha ,\beta )`$ is a central simple $`k`$-algebra. For our purposes we will assume that the subalgebras $`C(\alpha )=kx:x^n=\alpha `$ and $`C(\beta )=ky:y^n=\beta `$ are fields. Let $`\mathrm{\Sigma }`$ be an $`n`$-dimensional $`C(\alpha )`$-vector space with basis $`B=\{v_1,\mathrm{},v_n\}`$ and consider a dual basis $`B^{}=\{v_1^{},\mathrm{},v_n^{}\}`$ in $`\mathrm{\Sigma }^{}.`$ The algebra $`A_\omega (\alpha ,\beta )`$ can be embedded in $`M_n(C(\alpha ))`$ by assigning $$xX=xe_{1,1}+\omega xe_{2,2}+\mathrm{}+\omega ^{n1}xe_{n,n},yY=e_{1,2}+\mathrm{}+e_{n1,n}+\beta e_{n,1},$$ where $`e_{i,j}`$ denotes the elementary matrix in $`M_n(C(\alpha ))`$ with $`1`$ in the $`(i,j)`$-entry and zero elsewhere. The action of $`X`$ and $`Y`$ on the bases $`B`$ and $`B^{}`$ is: \| $`Xv_j=\omega ^{j1}v_jx`$ \| \| $`v_j^{}X=\omega ^{j1}xv_j^{}`$ \| \| --- \| --- \| --- \| \| $`Yv_j=\{\begin{array}{cc}\beta v_n\hfill & \text{if }j=1\hfill \\ v_{j1}\hfill & \text{if }j>1\hfill \end{array}`$ \| \| $`v_j^{}Y=\{\begin{array}{cc}\beta v_1^{}\hfill & \text{if }j=n\hfill \\ v_{j+1}^{}\hfill & \text{if }j<n\hfill \end{array}`$ \| (7) If either $`\alpha `$ or $`\beta `$ is equal to $`1`$, then the algebra $`A_\omega (\alpha ,\beta )`$ is isomorphic to $`M_n(k)`$. We will next describe the coideals of the $`C(\alpha )`$–coring $`\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$ corresponding to the intermediate extensions of $`kM_n(C(\alpha ))`$ given in the following diagram: For $`l,m,i=1,\mathrm{},n`$ set $`z_{l,m}^i=\alpha ^1x^{ni}v_l^{}_kv_mx^i`$. Observe that the set $`\{z_{l,m}^i:l,m,i=1,\mathrm{},n\}`$ is a basis of $`\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$ as a right $`C(\alpha )`$-vector space. We have that $`xz_{l,m}^i=z_{l,m}^{i1}x`$ for $`1in`$ with the convention $`z_{l,m}^0=z_{l,m}^n.`$ The comultiplication and counit of $`\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$ reads: $$\begin{array}{cc}\mathrm{\Delta }(z_{l,m}^i)\hfill & =_{j=1}^n(\alpha ^1x^{ni}v_l^{}_kv_j)_{C\left(\alpha \right)}(v_j^{}_kv_mx^i)\hfill \\ & =_{j=1}^n(\alpha ^1x^{ni}v_l^{}_kv_jx^i)_{C\left(\alpha \right)}(\alpha ^1x^{ni}v_j^{}_kv_mx^i)\hfill \\ & =_{j=1}^nz_{l,j}^i_{C\left(\alpha \right)}z_{j,m}^i,\hfill \\ ϵ(z_{l,m}^i)\hfill & =\delta _{l,m}.\hfill \end{array}$$ (8) Trivial extensions: The trivial extensions $`kkM_n(C(\alpha ))`$ and $`kM_n(C(\alpha ))M_n(C(\alpha ))`$ correspond to the coideals $`\{0\}`$ and $`Ker(ϵ)`$ respectively. Extension $`kC(\alpha )M_n(C(\alpha )):`$ We embed $`C(\alpha )`$ into $`M_n(C(\alpha ))`$ by mapping $`x`$ to $`X`$ and consider the $`C(\alpha )`$-coring $`\mathrm{\Sigma }^{}_{C\left(\alpha \right)}\mathrm{\Sigma }`$. The action of $`X`$ on $`B`$ and $`B^{}`$ gives the following relations in $`\mathrm{\Sigma }^{}_{C\left(\alpha \right)}\mathrm{\Sigma }:`$ $$\alpha ^1x^{ni}v_l^{}_{C\left(\alpha \right)}v_mx^i=v_l^{}_{C\left(\alpha \right)}v_m.$$ (9) The set $`\{v_l^{}_{C\left(\alpha \right)}v_m:l,m=1,\mathrm{},n\}`$ is a basis of $`\mathrm{\Sigma }^{}_{C\left(\alpha \right)}\mathrm{\Sigma }`$ as a right $`C(\alpha )`$-vector space. Set $`c_{l,m}=v_l^{}_{C\left(\alpha \right)}v_m`$. The bimodule structure on this coring is given by $`xc_{l,m}=\omega ^{ml}c_{l,m}x`$. The comultiplication and counit in $`\mathrm{\Sigma }^{}_{C\left(\alpha \right)}\mathrm{\Sigma }`$ are defined by: $$\mathrm{\Delta }(c_{l,m})=\underset{j=1}{\overset{n}{}}c_{l,j}_{C\left(\alpha \right)}c_{j,m},ϵ(c_{l,m})=\delta _{l,m}.$$ The coideal $`J_{C(\alpha )}`$ of $`\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$ corresponding to this extension is the right subspace generated by the set $$\{z_{l,m}^n\omega ^{i(ml)}z_{l,m}^i:l,m,i=1,\mathrm{},n1\}.$$ Observe that if we embed diagonally $`C(\alpha )`$ into $`M_n(C(\alpha ))`$, then $`\mathrm{\Sigma }`$ is centralized by $`C(\alpha )`$ and hence $`\mathrm{\Sigma }^{}_{C\left(\alpha \right)}\mathrm{\Sigma }`$ is indeed a coalgebra, the comatrix coalgebra over $`C(\alpha )`$ of order $`n`$. Hence the diagonal embedding of $`C(\alpha )`$ is not conjugated with the preceding one. Extension $`kC(\beta )M_n(C(\alpha )):`$ We embed $`C(\beta )`$ into $`M_n(C(\alpha ))`$ by mapping $`y`$ to $`Y`$ and consider the $`C(\alpha )`$-coring $`\mathrm{\Sigma }^{}_{C\left(\beta \right)}\mathrm{\Sigma }`$. Taking into account the action of $`Y`$ on $`B`$ and $`B^{}`$, the following relations in $`\mathrm{\Sigma }^{}_{C\left(\beta \right)}\mathrm{\Sigma }`$ are obtained: $$v_i^{}_{C\left(\beta \right)}v_j=\{\begin{array}{cc}\beta ^1v_{n(ji)+1}^{}_{C\left(\beta \right)}v_1\hfill & \text{if }i<j\hfill \\ v_{ij+1}^{}_{C\left(\beta \right)}v_1\hfill & \text{if }ij\hfill \end{array}$$ (10) A basis of $`\mathrm{\Sigma }^{}_{C\left(\beta \right)}\mathrm{\Sigma }`$ as a left $`C(\alpha )`$-vector space is: $$\{\alpha ^1x^{ni}v_l^{}_{C\left(\beta \right)}v_1x^i:l,i=1,\mathrm{},n\}.$$ Setting $`c_l^i=\alpha ^1x^{ni}v_l^{}_{C\left(\beta \right)}v_1x^i`$, the left action of $`C(\alpha )`$ on $`\mathrm{\Sigma }^{}_{C\left(\beta \right)}\mathrm{\Sigma }`$ reads as $`xc_l^i=c_l^{i1}x`$ with the convention $`c_l^0=c_l^n.`$ The comultiplication and counit of $`\mathrm{\Sigma }^{}_{C\left(\beta \right)}\mathrm{\Sigma }`$ is given by: $$\begin{array}{cc}\mathrm{\Delta }(c_l^i)\hfill & =_{j=1}^n(\alpha ^1x^{ni}v_l^{}_{C\left(\beta \right)}v_j)_{C\left(\alpha \right)}(v_j^{}_{C\left(\beta \right)}v_1x^i)\hfill \\ & =_{j=1}^n(\alpha ^1x^{ni}v_l^{}_{C\left(\beta \right)}v_jx^i)_{C\left(\alpha \right)}(\alpha ^1x^{ni}v_j^{}_{C\left(\beta \right)}v_1x^i)\hfill \\ & =_{j=1}^l(\alpha ^1x^{ni}v_l^{}_{C\left(\beta \right)}v_jx^i)_{C\left(\alpha \right)}(\alpha ^1x^{ni}v_j^{}_{C\left(\beta \right)}v_1x^i)\hfill \\ & +_{j=l+1}^n(\alpha ^1x^{ni}v_l^{}_{C\left(\beta \right)}v_jx^i)_{C\left(\alpha \right)}(\alpha ^1x^{ni}v_j^{}_{C\left(\beta \right)}v_1x^i)\hfill \\ & =_{j=1}^l(\alpha ^1x^{ni}v_{lj+1}^{}_{C\left(\beta \right)}v_1x^i)_{C\left(\alpha \right)}(\alpha ^1x^{ni}v_j^{}_{C\left(\beta \right)}v_1x^i)\hfill \\ & +\beta ^1_{j=l+1}^n(\alpha ^1x^{ni}v_{n(jl)+1}^{}_{C\left(\beta \right)}v_1x^i)_{C\left(\alpha \right)}(\alpha ^1x^{ni}v_j^{}_{C\left(\beta \right)}v_1x^i)\hfill \\ & =_{j=1}^lc_j^i_{C\left(\alpha \right)}c_{lj+1}^i+\beta ^1_{j=l+1}^nc_j^i_{C\left(\alpha \right)}c_{n(jl)+1}^i,\hfill \\ ϵ(c_l^i)\hfill & =\delta _{l,1}.\hfill \end{array}$$ The coideal $`J_{C(\beta )}`$ of $`\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$ corresponding to $`C(\beta )`$ is the right subspace generated by the following set: $$\begin{array}{c}\{z_{l,m}^i\beta ^1z_{n(ml)+1,1}^i:l,m,i=1,\mathrm{},n;l<m\}\{z_{l,m}^iz_{lm+1,1}^i:l,m,i=1,\mathrm{},n;lm\}.\hfill \end{array}$$ Extension $`kA_\omega (\alpha ,\beta )M_n(C(\alpha )):`$ Consider $`A_\omega (\alpha ,\beta )`$ as embedded into $`M_n(C(\alpha ))`$ by mapping $`x`$ to $`X`$ and $`y`$ to $`Y`$. We next describe the $`C(\alpha )`$-coring $`\mathrm{\Sigma }^{}_{A_\omega (\alpha ,\beta )}\mathrm{\Sigma }`$. Since $`C(\alpha )`$ and $`C(\beta )`$ are contained in $`A_\omega (\alpha ,\beta )`$, similar relations to (9) and (10) are obtained. The set $`\{v_i_{A_\omega (\alpha ,\beta )}v_1:i=1,\mathrm{},n\}`$ is a basis of $`\mathrm{\Sigma }^{}_{A_\omega (\alpha ,\beta )}\mathrm{\Sigma }`$ as a right $`C(\alpha )`$-vector space. Set $`c_i=v_i_{A_\omega (\alpha ,\beta )}v_1`$. The bimodule structure of $`\mathrm{\Sigma }^{}_{A_\omega (\alpha ,\beta )}\mathrm{\Sigma }`$ is $`xc_i=\omega ^{i+1}c_ix`$. The comultiplication and the counit of $`\mathrm{\Sigma }^{}_{A_\omega (\alpha ,\beta )}\mathrm{\Sigma }`$ are: $$\begin{array}{cc}\mathrm{\Delta }(c_i)\hfill & =_{l=1}^n(v_i^{}_{A_\omega (\alpha ,\beta )}v_l)_{C\left(\alpha \right)}(v_l^{}_{A_\omega (\alpha ,\beta )}v_1)\hfill \\ & =_{l=1}^i(v_{il+1}^{}_{A_\omega (\alpha ,\beta )}v_1)_{C\left(\alpha \right)}(v_l^{}_{A_\omega (\alpha ,\beta )}v_1)\hfill \\ & +\beta ^1_{l=i+1}^n(v_{n(li)+1}^{}_{A_\omega (\alpha ,\beta )}v_1)_{C\left(\alpha \right)}(v_l^{}_{A_\omega (\alpha ,\beta )}v_1),\hfill \\ & =_{l=1}^ic_l_{C\left(\alpha \right)}c_{il+1}+\beta ^1_{l=i+1}^nc_l_{C\left(\alpha \right)}c_{n(li)+1},\hfill \\ ϵ(c_i)\hfill & =\delta _{i,1}.\hfill \end{array}$$ The coideal $`J_{A_\omega (\alpha ,\beta )}`$ of $`\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$ that corresponds to this intermediate extension is the right subspace generated by the set: $$\begin{array}{c}\{z_{l,m}^n\omega ^{i(ml)}z_{l,m}^i:l,m,i=1,\mathrm{},n\}\{z_{l,m}^i\beta ^1z_{n(ml)+1,1}^i:l,m,i=1,\mathrm{},n;\hfill \\ l<m\}\{z_{l,m}^iz_{lm+1,1}^i:l,m,i=1,\mathrm{},n;lm\}.\hfill \end{array}$$ Extension $`kM_n(k)M_n(C(\alpha )):`$ Let $$X=e_{1,1}+\omega e_{2,2}+\mathrm{}+\omega ^{n1}e_{n,n},Y=e_{1,2}+\mathrm{}+e_{n1,n}+e_{n,1}.$$ Then $`X^n=1,Y^n=1`$ and $`YX=\omega XY`$ and the $`k`$-algebra generated by $`X`$ and $`Y`$ is $`M_n(k)`$. The action of $`X`$ and $`Y`$ on the basis $`B`$ and $`B^{}`$ is obtained from (7) for $`\beta =1`$. From these actions we get the relations: $$\begin{array}{cc}v_1^{}_{M_n\left(k\right)}v_1\hfill & =v_n^{}Y_{M_n\left(k\right)}v_1=v_n^{}_{M_n\left(k\right)}Yv_1=v_n^{}_{M_n\left(k\right)}v_n\hfill \\ & =v_{n1}^{}Y_{M_n\left(k\right)}v_n=v_{n1}^{}_{M_n\left(k\right)}Yv_n=v_{n1}^{}_{M_n\left(k\right)}v_{n1}\hfill \\ & =\mathrm{}.\hfill \\ & =v_2^{}_{M_n\left(k\right)}v_2.\hfill \\ v_i^{}_{M_n\left(k\right)}v_j\hfill & =\omega ^{j+1}v_i^{}_{M_n\left(k\right)}Xv_j=\omega ^{j+1}v_i^{}X_{M_n\left(k\right)}v_j\hfill \\ & =\omega ^{ij}v_i^{}_{M_n\left(k\right)}v_j.\hfill \end{array}$$ Then $`v_i^{}_{M_n\left(k\right)}v_j=0`$ for $`ij`$. Set $`c_i=\alpha ^1x^{ni}v_1^{}_{M_n\left(k\right)}v_1x^{ni}`$ for $`i=1,\mathrm{},n`$. Then the set $`\{c_i:i=1,\mathrm{},n\}`$ is a basis of $`\mathrm{\Sigma }^{}_{M_n\left(k\right)}\mathrm{\Sigma }`$ as a right $`C(\alpha )`$-vector space. The bimodule structure of this coring is given by $`xc_i=c_{i1}x`$ with the convention $`c_0=c_n`$. The comultiplication and counit of $`\mathrm{\Sigma }^{}_{M_n\left(k\right)}\mathrm{\Sigma }`$ takes the form: $$\begin{array}{cc}\mathrm{\Delta }(c_i)\hfill & =_{j=1}^n(\alpha x^{ni}v_1^{}_{M_n\left(k\right)}v_j)_{C\left(\alpha \right)}(v_j^{}_{M_n\left(k\right)}v_1)x^i\hfill \\ & =(\alpha x^{ni}v_1^{}_{M_n\left(k\right)}v_1)_{C\left(\alpha \right)}(v_1^{}_{M_n\left(k\right)}v_1)x^i\hfill \\ & =(\alpha x^{ni}v_1^{}_{M_n\left(k\right)}v_1x^i)_{C\left(\alpha \right)}(\alpha ^1x^{ni}v_1^{}_{M_n\left(k\right)}v_1)x^i\hfill \\ & =c_i_{C\left(\alpha \right)}c_i,\hfill \\ ϵ(c_i)\hfill & =1.\hfill \end{array}$$ This coring can also be obtained as the Sweedler coring associated to the extension $`kC(\alpha )`$. The coideal $`J_{M_n(k)}`$ of $`\mathrm{\Sigma }^{}_k\mathrm{\Sigma }`$ corresponding to this extension is the right subspace spanned by the set $$\{z_{l,m}^i:l,m,i=1,\mathrm{},n;lm\}\{z_{1,1}^iz_{l,l}^i:l,i=1,\mathrm{},n\}.$$ ## 3 Duality Let $`f:𝔇,g:𝔈`$ be surjective homomorphisms of $`A`$–corings. Then $`Kerf=Kerg`$ if and only if there exists an isomorphism of $`A`$–corings $`𝔇𝔈`$ making commute the diagram (11) Thus, every coideal $`J`$ of $``$ determines a class of surjective homomorphisms $`𝔇`$ having $`J`$ as their common kernel or, alternatively, the morphisms in each class are connected by commutative triangles as in (11). From a formal point of view, corings over $`A`$ are dual to $`A`$–rings, being these last understood to be morphisms of rings $`AU`$. The definition of a homomorphism of $`A`$–rings is obvious, and we will conceive an $`A`$–subring of a given $`A`$–ring $`AE`$ as an isomorphism class of injective homomorphisms of $`A`$–rings $`UE`$. Obviously, every $`A`$–subring of $`E`$ may be represented by an inclusion $`UE`$. We can thus consider the set $`\mathrm{𝖲𝗎𝖻𝗋𝗂𝗇𝗀𝗌}(E)`$ of $`A`$–subrings of $`E`$. One of the possible concrete dual correspondences from $`A`$–corings to $`A`$–rings goes as follows: if $``$ is an $`A`$–coring, then $`{}_{}{}^{}ϵ_{}^{}:A{}_{}{}^{}_{}^{op}`$ is an $`A`$–ring (see \[14, Proposition 3.2\]), and under this mapping, homomorphisms of $`A`$–corings give homomorphisms of $`A`$–rings. In particular, if $`J\mathrm{𝖢𝗈𝗂𝖽𝖾𝖺𝗅𝗌}()`$ and $`𝔇=/J`$ is the corresponding canonical projection, then we have an injective homomorphism of $`A`$–rings $`{}_{}{}^{}𝔇_{}^{op}{}_{}{}^{}_{}^{op}`$ (see again \[14, Proposition 3.2\]). If $`=\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$ is a comatrix $`A`$–coring, then, by \[6, Proposition 2.1\], we have an injective homomorphism of $`A`$–rings $`{}_{}{}^{}𝔇_{}^{op}{}_{}{}^{}_{}^{op}\mathrm{End}({}_{B}{}^{}\mathrm{\Sigma })`$. The corresponding $`A`$–subring of $`\mathrm{End}({}_{B}{}^{}\mathrm{\Sigma })`$ will be denoted by $`^{}(J)`$. Conversely, given an $`A`$–subring $`U\mathrm{End}({}_{B}{}^{}\mathrm{\Sigma })`$, then $`\mathrm{End}(\mathrm{\Sigma }_U)`$ is independent on the representative $`U`$ of the $`A`$–subring. We define the coideal $`𝒥^{}(U)`$ of $`\mathrm{\Sigma }^{}_B\mathrm{\Sigma }`$ as the kernel of the homomorphism of $`A`$–corings $`\mathrm{\Sigma }^{}_B\mathrm{\Sigma }\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_U)}\mathrm{\Sigma }`$ induced by the ring homomorphism $`B\mathrm{End}(\mathrm{\Sigma }_U)`$. These considerations lead to a Galois connection for $`E=\mathrm{End}({}_{B}{}^{}\mathrm{\Sigma })`$: (12) The following proposition collects some of its relevant properties. A right comodule $`\mathrm{\Sigma }`$ over an $`A`$–coring $`𝔇`$ is said to be loyal if the canonical map induced by the inclusion $`\mathrm{End}(\mathrm{\Sigma }_𝔇)\mathrm{End}({}_{{}_{}{}^{}𝔇}{}^{}\mathrm{\Sigma })`$ is a bijection. By \[3, 19.2,19.3\], if $`{}_{A}{}^{}𝔇`$ is locally projective, then every right $`𝔇`$–comodule is loyal. ###### Proposition 3.1. The mappings defined in (12) enjoy the following properties. 1. $`^{}𝒥^{}(U)=\mathrm{End}({}_{\mathrm{End}(\mathrm{\Sigma }_U)}{}^{}\mathrm{\Sigma })`$ for every $`U\mathrm{𝖲𝗎𝖻𝗋𝗂𝗇𝗀𝗌}(E)`$ and, thus, $`U^{}𝒥^{}(U)`$. 2. $`J𝒥^{}^{}(J)`$ for every $`J\mathrm{𝖢𝗈𝗂𝖽𝖾𝖺𝗅𝗌}(\mathrm{\Sigma }^{}_B\mathrm{\Sigma })`$ such that $`\mathrm{\Sigma }_{/J}`$ is loyal. 3. The maps $`𝒥^{}`$ and $`^{}`$ establish a bijection between the set of $`A`$–subrings $`U`$ of $`\mathrm{End}({}_{B}{}^{}\mathrm{\Sigma })`$ such that $`\mathrm{\Sigma }_U`$ is faithful and balanced, and the set of coideals $`J`$ of $``$ such that $`(/J,\mathrm{\Sigma })`$ is Galois and $`\mathrm{\Sigma }_{/J}`$ is loyal. ###### Proof. (1) Given an $`A`$–subring $`U\mathrm{End}(\mathrm{\Sigma }_B)`$, $`𝒥^{}(U)`$ is defined as the kernel of the surjective homomorphism of $`A`$–corings $`\mathrm{\Sigma }^{}_B\mathrm{\Sigma }\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_U)}\mathrm{\Sigma }`$. The commutative diagram of injective homomorphisms of $`A`$–rings deduced from \[6, Proposition 2.1\], shows that $`^{}𝒥^{}(U)=\mathrm{End}({}_{\mathrm{End}(\mathrm{\Sigma }_U)}{}^{}\mathrm{\Sigma })`$. (2) With $`p:𝔇=/J`$ the canonical projection, consider the commutative diagram $$\text{},$$ (13) where $`f`$ is induced by the morphism $`B\mathrm{End}(\mathrm{\Sigma }_𝔇)`$. Since $`\mathrm{\Sigma }_𝔇`$ is loyal, we get from the diagram that $`J\mathrm{ker}(f)=𝒥^{}^{}(J)`$. (3) Let $`UE=\mathrm{End}({}_{B}{}^{}\mathrm{\Sigma })`$ be an $`A`$–subring. By definition, $`𝒥^{}(U)`$ is such that $`/𝒥^{}(U)\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_U)}\mathrm{\Sigma }`$. By \[6, Lemma 3.9\], $`(\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_U)}\mathrm{\Sigma },\mathrm{\Sigma })`$ is Galois. This means that the canonical map is an isomorphism. But this map is nothing but the one induced by the ring extension $$\mathrm{End}(\mathrm{\Sigma }_{\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_U)}\mathrm{\Sigma }})\mathrm{End}({}_{{}_{}{}^{}(\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_U)}\mathrm{\Sigma })}{}^{}\mathrm{\Sigma })=\mathrm{End}(\mathrm{\Sigma }_{\mathrm{End}(\mathrm{\Sigma }_{\mathrm{End}(\mathrm{\Sigma }_U)})})=\mathrm{End}(\mathrm{\Sigma }_U),$$ which proves that $`\mathrm{\Sigma }`$ is a loyal right $`\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_U)}\mathrm{\Sigma }`$–comodule. Part (1) gives obviously $`^{}𝒥^{}(U)=U`$. Conversely, let $`J`$ be a coideal of $``$ such that $`(/J,\mathrm{\Sigma })`$ is Galois and $`\mathrm{\Sigma }_{/J}`$ is Galois. Put $`𝔇=/J`$. From the triangle (13), and the fact that $`\mathrm{\Sigma }_𝔇`$ is loyal, we compute the $`A`$–subring $`^{}(J)`$ of $`\mathrm{End}({}_{R}{}^{}\mathrm{\Sigma })`$ as $${}_{}{}^{}(𝔇)_{}^{op}{}_{}{}^{}(\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_𝔇)}\mathrm{\Sigma })_{}^{op}{}_{}{}^{}(\mathrm{\Sigma }^{}_{\mathrm{End}(\mathrm{\Sigma }_{{}_{}{}^{}𝔇_{}^{op}})}\mathrm{\Sigma })_{}^{op}\mathrm{End}({}_{\mathrm{End}(\mathrm{\Sigma }_{{}_{}{}^{}𝔇_{}^{op}})}{}^{}\mathrm{\Sigma })$$ (14) This means that $`^{}(J)`$, defined as $`{}_{}{}^{}𝔇_{}^{op}`$, is such that $`\mathrm{\Sigma }_{^{}(J)}`$ is faithful and balanced. The diagram (13) gives in addition that $`𝒥^{}^{}(J)=J`$. ∎ We are now in position to prove our version for simple artinian rings of the Jacobson-Bourbaki theorem. ###### Theorem 3.2. Let $`\mathrm{\Sigma }`$ be a finitely generated projective right module over a Quasi-Frobenius ring $`A`$ and let $`S=\mathrm{End}(\mathrm{\Sigma }_A)`$. Let $`B`$ be a subring of $`S`$. Denote by $``$ the set of intermediate simple artinian subrings $`BCS`$ such that $`{}_{C}{}^{}\mathrm{\Sigma }`$ is finitely generated. Let $`𝒟`$ denote the set of all $`A`$-subrings $`U`$ of $`\mathrm{End}({}_{B}{}^{}\mathrm{\Sigma })`$ such that $`U_A`$ is finitely generated and projective, and $`\mathrm{\Sigma }_U`$ is semisimple and isotypic. The maps $$\begin{array}{c}^{}():𝒟,U\mathrm{End}(\mathrm{\Sigma }_U),\hfill \\ 𝒥^{}():𝒟,C\mathrm{End}({}_{C}{}^{}\mathrm{\Sigma }).\hfill \end{array}$$ establish a bijective correspondence. This correspondence is dual to that of Theorem 2.4. Moreover, if $`A`$ is in addition local, then two intermediate simple artinian rings $`C`$ and $`C^{}`$ are conjugated if and only if $`\mathrm{End}({}_{C}{}^{}A)`$ and $`\mathrm{End}({}_{C^{}}{}^{}A)`$ are isomorphic as $`A`$–rings. ###### Proof. If $`C`$ is simple artinian and $`{}_{C}{}^{}\mathrm{\Sigma }`$ is finitely generated, then $`U=\mathrm{End}({}_{C}{}^{}\mathrm{\Sigma })`$ is a simple artinian $`A`$–ring. Thus, $`\mathrm{\Sigma }_U`$ is semisimple isotypic. The comatrix $`A`$–coring $`\mathrm{\Sigma }^{}_C\mathrm{\Sigma }`$ is then finitely generated and projective as a left $`A`$–module. The ring isomorphism $`U^{op}{}_{}{}^{}(\mathrm{\Sigma }^{}_C\mathrm{\Sigma })`$ given in \[6, Proposition 2.1\] is an isomorphism of $`A`$–bimodules. In particular, $`U_A`$ becomes a finitely generated module. Obviously, $`{}_{C}{}^{}\mathrm{\Sigma }`$ is faithful and balanced. Conversely, assume that $`\mathrm{\Sigma }_U`$ is semisimple isotypic for an $`A`$–subring $`U\mathrm{End}({}_{B}{}^{}\mathrm{\Sigma })`$ such that $`U_A`$ is finitely generated. Then $`C=\mathrm{End}(\mathrm{\Sigma }_U)`$ is a simple artinian subring of $`S`$, because $`\mathrm{\Sigma }_U`$ is clearly finitely generated. To prove that $`\mathrm{\Sigma }_U`$ is faithful and balanced, consider that the comatrix $`A`$–coring structure on $`U^{}_UU`$ induces, via the isomorphism $`U^{}_UUU^{}`$ an $`A`$–coring structure on this last $`A`$–bimodule \[6, Example 2.4\]. If $`\{u_\alpha ^{},u_\alpha \}_{\alpha =1}^n`$ is a finite dual basis for $`U_A`$, the comultiplication and counit are given explicitly by $$\mathrm{\Delta }:U^{}U^{}_AU^{},\phi \underset{\alpha }{}\phi u_\alpha _Au_\alpha ^{},$$ $$ϵ:U^{}A,\phi \phi (1)$$ The canonical map $`U{}_{}{}^{}(U^{})`$ is then an anti-isomorphism of rings. We may then identify the categories of $`^U^{}`$ and $`{}_{U^{op}}{}^{}=_U`$ (see e.g. \[3, 19.6\]). In fact, an explicit isomorphism of categories is given as follows. For each element $`m`$ in a right $`U`$–module $`M`$, the equality $`mu=_\alpha mu_\alpha u_\alpha ^{}(u)`$ for $`uU`$ says that $`\{mu_\alpha ,u_\alpha ^{}\}`$ is a set of right rational parameters in the sense of . Thus \[7, Corollary 4.7\] gives the isomorphism of categories $`_U=^U^{}`$, where the right $`U^{}`$–comodule structure on $`M_U`$ is given by $`\rho _M(m)=_\alpha mu_\alpha _Au_\alpha ^{}`$. In particular, the (finitely generated) simple isotypic right $`U`$–module may be considered as a right $`U^{}`$–comodule, and we have a canonical map given by $$\mathrm{𝖼𝖺𝗇}:\mathrm{\Sigma }^{}_C\mathrm{\Sigma }U^{},\phi _Cx\underset{\alpha }{}\phi (xu_\alpha )u_\alpha ^{}$$ (15) Now, for every $`0uU`$, let $`x\mathrm{\Sigma }`$ such that $`xu0`$. Then, for $`\phi \mathrm{\Sigma }^{}`$ with $`\phi (xu)0`$, we have $$\mathrm{𝖼𝖺𝗇}(\phi _Cx)(u)=\underset{\alpha }{}\phi (xu_\alpha )u_\alpha ^{}(u)=\underset{\alpha }{}\phi (xu_\alpha u_\alpha ^{}(u))=\phi (xu)0,$$ which implies, being $`A`$ Quasi-Frobenius and $`U_A`$ finitely generated, that $`\mathrm{𝖼𝖺𝗇}`$ is surjective. Finally, since $`\mathrm{\Sigma }_U^{}`$ is semisimple, we get from Corollary 1.6 that $`\mathrm{𝖼𝖺𝗇}`$ an isomorphism. Therefore, we have an isomorphism which turns out to be, by using \[6, equation (3)\], the canonical homomorphism $`U\mathrm{End}({}_{C}{}^{}\mathrm{\Sigma })`$. In this way, $`\mathrm{\Sigma }_U`$ is faithful and balanced and $`{}_{C}{}^{}\mathrm{\Sigma }`$ must be finitely generated. ∎ ###### Remark 3.3. If the base ring $`A`$ is simple artinian then Theorem 3.2 takes a simpler form. Thus, for a simple artinian ring $`B`$, the maps $$\begin{array}{c}^{}():𝒟,UEnd(\mathrm{\Sigma }_U),\hfill \\ 𝒥^{}():𝒟,CEnd({}_{C}{}^{}\mathrm{\Sigma }).\hfill \end{array}$$ establish a bijective correspondence between the set $`𝒟`$ of simple artinian $`A`$-subrings $`U`$ of $`End({}_{B}{}^{}\mathrm{\Sigma })`$ such that $`U_A`$ is finitely generated and the set $``$ of intermediate simple artinian subrings $`BCS`$ such that $`{}_{C}{}^{}S`$ is finitely generated. If $`\mathrm{\Sigma }=A`$, and $`BA`$ is a subring, then for every $`A`$–subring $`U\mathrm{End}({}_{B}{}^{}A)`$, we have that $`^{}(U)=\mathrm{End}(U_A)=\{cA:u(ca)=cu(a),uU,aA\}`$. If, moreover, $`A`$ is a division ring, then we deduce from Theorem 3.2 the following version of the Jacobson-Bourbaki theorem (c.f. \[14, §4\]). ###### Corollary 3.4. Let $`BA`$ be division rings. The maps $`^{}`$ and $`𝒥^{}`$ establish a bijective correspondence between the intermediate division rings $`BCA`$ such that $`{}_{C}{}^{}A`$ is finite dimensional and the $`A`$–subrings $`U\mathrm{End}({}_{B}{}^{}A)`$ such that $`U_A`$ is finite dimensional. Moreover, two intermediate division subrings $`C`$ and $`C^{}`$ are conjugated if and only if $`\mathrm{End}({}_{C}{}^{}A)`$ and $`\mathrm{End}({}_{C^{}}{}^{}A)`$ are isomorphic $`A`$–rings. ###### Proof. The pertinent remark here is that $`A_U`$ is always a simple module (since $`A_A`$ is simple). ∎ ## Appendix Throughout this section all corings are $`/`$-corings and all $``$-bimodules are assumed to be centralized by the field of real numbers $``$. We know by the Structure Theorem of simple and cosemisimple corings that any such a $``$-coring $``$ is of the form $`\mathrm{\Sigma }^{}_D\mathrm{\Sigma }`$ where $`\mathrm{\Sigma }`$ is a finite dimensional complex vector space and $`D`$ is a $``$-division algebra embedded in $`\mathrm{End}(\mathrm{\Sigma }_{})M_n()`$ where $`n=dim(\mathrm{\Sigma }_{}).`$ Furthermore, two corings $`\mathrm{\Sigma }^{}_D\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}_E\mathrm{\Sigma }`$ are isomorphic if and only if there is an invertible $`u\mathrm{End}(\mathrm{\Sigma }_{})`$ such that $`uEu^1=D`$, i.e., $`E`$ and $`D`$ are conjugated in $`\mathrm{End}(\mathrm{\Sigma }_{})`$. By Fröbenius Theorem, $`D=,`$ or $``$. We study the possible ways of embedding these division algebras in $`M_n()`$. As a consequence of the proof of the Skolem-Noether Theorem (\[10, Theorem 4.9\]) the non conjugated ways of embedding $`D`$ in $`M_n()`$ are in bijective correspondence with the simple left $`D_{}M_n()`$-modules. 1. Case $`D=`$. Since $`M_n()`$ is an $``$-algebra, the only way of embedding $``$ in $`M_n()`$ is the obvious one. The comultiplication and counit of the coring $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ is described in (8). If $`n=1`$, then this coring is isomorphic to the coring $`[/2]`$ associated to the canonical $`/2`$–grading on $``$. As a right $``$–vector space, $`[/2]`$ is free over the basis $`/2=\{[0],[1]\}`$. Its left $``$-vector space structure is determined by the rules $`i[0]=[1]i`$, $`i[1]=[0]i`$. In this coring, $`[0]`$ and $`[1]`$ are group-like elements. For $`n>1`$, the coring $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ is isomorphic with the tensor product coring (see \[8, Proposition 1.5\]) $`[/2]_{}M^c(,n)`$, where $`M^c(,n)`$ is the comatrix $``$–coalgebra. An explicit isomorphism sends $`z_{m,l}^0`$ onto $`[0]_{}x_{m,l}`$ and $`z_{m,l}^1`$ onto $`[1]_{}x_{m,l}`$, where $`x_{m,l}`$ denotes the matrix with $`1`$ in the component $`(l,m)`$ and $`0`$ elsewhere. 2. Case $`D=`$. Since $`_{}M_n()M_n()M_n()`$, there are two ways (for $`n>0`$) of embedding $``$ into $`M_n()`$. Let $`e_{p,q}`$ denote the elementary matrix in $`M_n()`$ with $`1`$ in the $`(p,q)`$-entry and zero elsewhere. The two non conjugate embeddings are represented by the one sending $`i`$ to $`i(_{l=1}^ne_{l,l})`$ and the one sending $`i`$ to $`\overline{i}=_{l=1}^ne_{l,nl+1}.`$ In the first case, $`\mathrm{\Sigma }`$ is centralized by $``$ and therefore $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ is a $``$-coalgebra. Thus $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ is isomorphic to the comatrix coalgebra of order $`n`$. Let us study the second case. The action of $`\overline{i}`$ on $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$ is: $$\overline{i}v_j=v_{nj+1}i,v_j^{}\overline{i}=iv_{nj+1}^{}.$$ Then the bimodule structure on $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ is given by: $$\begin{array}{cc}i(v_p^{}_{}v_q)\hfill & =iv_p^{}_{}v_q=v_{np+1}^{}\overline{i}_{}v_q=v_{np+1}^{}_{}\overline{i}v_q\hfill \\ & =v_{np+1}^{}_{}v_{nq+1}i=(v_{np+1}^{}_{}v_{nq+1})i.\hfill \end{array}$$ The comultiplication and counit of $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ is defined by: $$\mathrm{\Delta }(v_p^{}v_q)=\underset{l=1}{\overset{n}{}}(v_p^{}v_l)(v_l^{}v_q),ϵ(v_p^{}v_q)=\delta _{p,q}.$$ This coring can be described as the right $``$-vector space $`=_{p,q=1}^nv_{p,q}`$ with bimodule structure $`iv_{p,q}=v_{np+1,nq+1}i`$ and with comultiplication and counit given by: $$\mathrm{\Delta }(v_{p,q})=\underset{l=1}{\overset{n}{}}v_{p,l}_{}v_{l,q},ϵ(v_{p,q})=\delta _{p,q}.$$ 3. Case $`D=`$. Since $``$ is a central simple $``$-algebra and $`M_n()`$ is simple, $`_{}M_n()`$ is simple. Hence all embeddings of $``$ in $`M_n()`$ are conjugate. Let us observe that if $``$ embeds in $`M_n()`$, then $`n`$ is even. By the Double Centralizer Theorem we would have $`M_n()_{}C()`$, where $`C()`$ denotes the centralizer of $``$ in $`M_n()`$. Comparing real dimensions, $`4`$ divides $`2n^2`$. Let $`\overline{i},\overline{j}`$ be the generators of $``$. Consider the following embedding of $``$ in $`M_2():`$ $$\overline{i}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\overline{j}\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right).$$ Since $`n`$ is even, the above embedding gives an embedding of $``$ in $`M_n()`$ by placing each block repeatedly in the main diagonal. Any other embedding of $``$ in $`M_n()`$ is conjugated to this one. We next describe the coring $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$. We first assume that $`n=2`$. Observe that this case is a particular case of our example in the second section by taking $`k=`$ and $`\alpha =\beta =1`$. The bimodule structure on $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ is the following: $$i(v_1^{}v_1)=(v_1^{}v_1)i,i(v_2^{}v_1)=(v_2^{}v_1)i.$$ The comultiplication and counit of $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ read as: $$\begin{array}{cc}\mathrm{\Delta }(v_1^{}v_1)=(v_1^{}v_1)(v_1^{}v_1)(v_2^{}v_1)(v_2^{}v_1),\hfill & ϵ(v_1^{}v_1)=1,\hfill \\ \mathrm{\Delta }(v_2^{}v_1)=(v_2^{}v_1)(v_1^{}v_1)+(v_1^{}v_1)(v_2^{}v_1),\hfill & ϵ(v_2^{}v_1)=0.\hfill \end{array}$$ This coring is precisely the trigonometric coring. We now discuss the general case $`dim(\mathrm{\Sigma }_{})=n=2m`$. Let us recall that $``$ is embedded in $`M_n()`$ in the following way: $$\begin{array}{cc}\overline{i}_{i=1}^me_{2l1,2l}_{l=1}^me_{2l,2l1},\overline{j}i(_{l=1}^n(1)^{l+1}e_{l,l}).\hfill & \end{array}$$ Through this embedding the action of $``$ on $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$ is: \| $`\overline{i}v_q=\{\begin{array}{cc}v_{q1}\hfill & \text{if q is even}\hfill \\ v_{q+1}\hfill & \text{if q is odd}\hfill \end{array}`$ \| \| $`\overline{j}v_q=(1)^{q+1}v_qi`$ \| \| --- \| --- \| --- \| \| $`v_q^{}\overline{i}=\{\begin{array}{cc}v_{q1}^{}\hfill & \text{if q is even}\hfill \\ v_{q+1}^{}\hfill & \text{if q is odd}\hfill \end{array}`$ \| \| $`v_q^{}\overline{j}=(1)^{q+1}iv_q^{}`$ \| These actions give rise to the following relations in $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$. Let $`p,q\{1,2,\mathrm{},m\}`$. Then: $$\begin{array}{c}v_{2p}^{}v_{2q}=v_{2p}^{}\overline{i}v_{2q1}=v_{2p}^{}\overline{i}v_{2q1}=v_{2p1}^{}v_{2q1},\hfill \\ v_{2p1}^{}v_{2q}=v_{2p1}^{}\overline{i}v_{2q}=v_{2p1}^{}\overline{i}v_{2q1}=v_{2p}^{}v_{2q1}.\hfill \end{array}$$ The set $`\{v_{2p}v_l\|p=1,\mathrm{},m;l=1,\mathrm{},n\}`$ is a basis of $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ as a right $``$-vector space. The left $``$-action on this coring is: $$\begin{array}{cc}i(v_{2p}^{}v_{2q})\hfill & =iv_{2p}^{}v_{2q}=(1)^{2p+1}v_{2p}^{}\overline{j}v_{2q}=(1)^{2p+1}v_{2p}^{}\overline{j}v_{2q}\hfill \\ & =v_{2p}^{}v_{2q}i=(v_{2p}^{}v_{2q})i,\hfill \\ i(v_{2p}^{}v_{2q1})\hfill & =iv_{2p}^{}v_{2q1}=(1)^{2p+1}v_{2p}^{}\overline{j}v_{2q1}=(1)^{2p+1}v_{2p}^{}\overline{j}v_{2q1}\hfill \\ & =v_{2p}^{}v_{2q}i=(v_{2p}^{}v_{2q1})i.\hfill \end{array}$$ The comultiplication and counit of $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ is: $$\begin{array}{cc}\mathrm{\Delta }(v_{2p}v_{2q})\hfill & =_{l=1}^m(v_{2p}^{}v_{2l})_{}(v_{2l}^{}v_{2q})+_{l=1}^m(v_{2p}^{}v_{2l1})_{}(v_{2l1}^{}v_{2q})\hfill \\ & =_{l=1}^m(v_{2p}^{}v_{2l})_{}(v_{2l}^{}v_{2q})_{l=1}^m(v_{2p}^{}v_{2l1})_{}(v_{2l}^{}v_{2q1})\hfill \\ ϵ(v_{2p}v_{2q})\hfill & =\delta _{p,q}\hfill \\ \mathrm{\Delta }(v_{2p}v_{2q1})\hfill & =_{l=1}^m(v_{2p}^{}v_{2l})_{}(v_{2l}^{}v_{2q1})+_{l=1}^m(v_{2p}^{}v_{2l1})_{}(v_{2l1}^{}v_{2q1})\hfill \\ & =_{l=1}^m(v_{2p}^{}v_{2l})_{}(v_{2l}^{}v_{2q})+_{l=1}^m(v_{2p}^{}v_{2l1})_{}(v_{2l}^{}v_{2q})\hfill \\ ϵ(v_{2p}v_{2q1})\hfill & =0.\hfill \end{array}$$ Let $`𝔗=cs`$ denote the trigonometric coring and let $`M^c(,m)`$ be the $``$-comatrix coalgebra of order $`m`$. Then $`𝔗_{}M^c(,m)`$ becomes a $``$-coring in the natural way . It may be verified that the map from $`𝔗_{}M^c(,m)`$ to $`\mathrm{\Sigma }^{}_{}\mathrm{\Sigma }`$ defined by $$c_{}x_{pq}v_{2p}^{}v_{2q},s_{}x_{pq}v_{2p}^{}v_{2q1},$$ is an isomorphism of corings.
warning/0506/cond-mat0506417.html	ar5iv	text	# Hybrid Burnett Equations. A New Method of Stabilizing ## I Introduction The practical interest of extending the Navier-Stokes equations to smaller length scales comes from the need to model re-entrance of space vehicles and in later times from nano scale technology. In this parameter region there is in general the Boltzmann equation. As the collision integral is very heavy to calculate numerically, the Boltzmann equation is usually applied for fairly simple or idealized problems, see Cercignani Cercignanis lilla bok . An other method is that of DSMC due to Bird Bird . But in the collision dominated region also DSMC is very demanding computationally. Hence special methods are valuable when the Knudsen number is too large for the Navier-Stokes equations to apply but still small enough to allow for an expansion. A fruitful technique in this region is the asymptotic method, originating with Hilbert and Grad, see Sone Sone . The Navier-Stokes equations were derived from the Boltzmann equation by the Chapman-Enskog method, see Chapman & Cowling CC . They are to first order in the mean free path. The corresponding equations to second order in the mean free path were derived by Burnett Burnett , see also CC . However, the Burnett equations were proven by Bobylev Bobylev to have a nonphysical instability. See also the review by Agarwal et al Agarwal and the paper by Uribe et al. Uribe . For the Burnett equations the trivial state of rest is thus unstable for perturbations of a wavelength of the order of the mean free path and shorter. The Chapman-Enskog method is an expansion in the Knudsen number $`Kn=l/L`$, where $`l`$ is the mean free path and $`L`$ a characteristic length. Wavelengths of the order of the mean free path and larger correspond to an effective Knudsen number of the order $`1`$ and larger. So there is no contradiction in the fact that the Burnett equations are unstable for short wavelengths. The physical content of the Burnett equations is for solutions with a characteristic length which is large enough. But nevertheless the equations should be wellbehaved for all length scales. This is necessary mathematically as well as numerically. As pointed out by Agarwal et al Agarwal Burnett’s original expression for the viscous pressure tensor contains the time derivative of the traceless rate of deformation tensor. Chapman & Cowling replace the time derivative by spatial derivatives using the equations at the Euler level. Agarwal et al give the name conventional Burnett equations to the resulting equations. There is a similar replacement for the time derivative of the temperature gradient. This method of replacement using the zero order equations is part of the procedure in the Chapman-Enskog method to obtain the Navier-Stokes equations from the Boltzmann equation to first order of the Knudsen number. Agarwal et al also raise the question if the instability of the conventional Burnett equations is caused by this replacement, so that possibly the original Burnett equations are linearly stable. In the present paper we show that the original Burnett equations have an unphysical singularity. In Jin Jin and Slemrod introduce $`𝐏,𝐪`$ as new independent fields and propose 13 equations, first order in time and second order in space. See also the two master theses by Svärd Svard and by Strömgren Stromgren and the paper by Söderholm Soderholm based on essentially the same idea but with equations which are first order in time as well as space. In the paper Soderholm these equations are applied numerically for nonlinear sound waves. In the present paper we introduce a two parameter partial replacement of the mentioned time derivatives and show that the parameters can be chosen to give linear stability. We then choose the parameters such that the resulting equations are as close to the original and conventional equations as possible. The resulting equations are more similar to the already studied Burnett equations than those in Jin and Soderholm and it seems that it should be simpler to reduce a numerical scheme for the equations to one for the Navier-Stokes equations. ## II The Burnett Equations The general equations of balance are $`{\displaystyle \frac{D\rho }{Dt}}+\rho 𝐯`$ $`=`$ $`0,`$ (1) $`\rho {\displaystyle \frac{D𝐯}{Dt}}`$ $`=`$ $`𝐏,`$ (2) $`\rho {\displaystyle \frac{3k_B}{2m}}{\displaystyle \frac{DT}{Dt}}`$ $`=`$ $`𝐏:𝐯𝐪.`$ (3) $`𝐏`$ is the pressure tensor and $`𝐪`$ the heat current. We are using dyadic notation. Let us first write down the original Burnett expression for the pressure tensor, see CC ($`p=k_B\rho T/m`$) $`𝐏_o`$ $`=`$ $`p\mathrm{𝟏}2\mu 𝐒+\varpi _1{\displaystyle \frac{\mu ^2}{p}}(𝐯)𝐒`$ $`+\varpi _2{\displaystyle \frac{\mu ^2}{p}}[{\displaystyle \frac{D𝐒}{Dt}}2𝐒(𝐯)]`$ $`+\varpi _3{\displaystyle \frac{\mu ^2}{\rho T}}T+\varpi _4{\displaystyle \frac{\mu ^2}{\rho pT}}pT`$ $`+\varpi _5{\displaystyle \frac{\mu ^2}{\rho T^2}}TT+\varpi _6{\displaystyle \frac{\mu ^2}{p}}𝐒𝐒.`$ Here, $`(𝐯)_{ij}`$ $`=`$ $`{\displaystyle \frac{}{x_i}}v_j=v_{j,i},(T)_{ij}=T_{,ij}`$ $`𝐒`$ $`=`$ $`𝐯={\displaystyle \frac{1}{2}}[𝐯+(𝐯)^T]{\displaystyle \frac{1}{3}}(𝐯)\mathrm{𝟏}.`$ $`\mathrm{𝟏}`$ is the unit tensor, $`\mathrm{}`$ means the symmetric traceless part. All other quantities have their usual meaning. The original Burnett expression for the heat current is $`𝐪_o`$ $`=`$ $`\kappa T+\theta _1{\displaystyle \frac{\mu ^2}{\rho T}}(𝐯)T`$ $`+\theta _2{\displaystyle \frac{\mu ^2}{\rho T}}[{\displaystyle \frac{D(T)}{Dt}}(𝐯)T]`$ $`+\theta _3{\displaystyle \frac{\mu ^2}{\rho p}}𝐒p+\theta _4{\displaystyle \frac{\mu ^2}{\rho }}𝐒+\theta _5{\displaystyle \frac{3\mu ^2}{\rho T}}𝐒T.`$ In Chapman & Cowling CC the time derivatives $`D/Dt`$ are replaced by spatial derivatives using the zero order equations. The resultting expressions are denoted $`D_0/Dt.`$ To start with we have $$\frac{D}{Dt}(T)=\frac{DT}{Dt}(𝐯)T.$$ The energy equation to zero order is $$\frac{3}{2}\frac{D_0T}{Dt}=T(𝐯).$$ Hence, $`{\displaystyle \frac{D_0}{Dt}}(T)`$ (6) $`={\displaystyle \frac{2}{3}}T(𝐯){\displaystyle \frac{2}{3}}(𝐯)T`$ $``$ $`(𝐯)T.`$ Similarly we have $`{\displaystyle \frac{D𝐯}{Dt}}`$ $`=`$ $`{\displaystyle \frac{D}{Dt}}𝐯+(𝐯)^2,`$ $`{\displaystyle \frac{D𝐯}{Dt}}`$ $`=`$ $`{\displaystyle \frac{D𝐒}{Dt}}+(𝐯)^2.`$ (7) The zero order momentum equation is $$\frac{D_0𝐯}{Dt}=\frac{1}{\rho }p.$$ Hence, $$\frac{D_0𝐒}{Dt}=(\frac{1}{\rho }p)𝐯)^2.$$ (8) Explicitly, $`{\displaystyle \frac{D_0𝐒}{Dt}}+(𝐯)^2`$ (9) $`={\displaystyle \frac{p}{\rho T}}{\displaystyle \frac{1}{\rho }}T\rho {\displaystyle \frac{T}{\rho ^2}}\rho \rho `$ $`+`$ $`{\displaystyle \frac{T}{\rho }}\rho +T.`$ Interpreting the time derivatives in (II) and (II ) according to (8) and (9) we find the conventional expressions for the pressure tensor and heat current. Let us denote them $`𝐏_𝐜,𝐪_𝐜.`$ The equations of balance (1) - (3) then give the conventional Burnett equations. ## III Stability and Replacing Time derivatives by space derivatives We now make the replacements $`{\displaystyle \frac{D𝐒}{Dt}}`$ $``$ $`(1\alpha ){\displaystyle \frac{D𝐒}{Dt}}+\alpha {\displaystyle \frac{D_0𝐒}{Dt}},`$ (10) $`{\displaystyle \frac{D}{Dt}}(T)`$ $``$ $`(1\beta ){\displaystyle \frac{D}{Dt}}(T)+\beta {\displaystyle \frac{D_0}{Dt}}(T),`$ (11) where $`\alpha ,\beta `$ are coefficients for which we later shall obtain bounds. The choice $`\alpha =\beta =0`$ gives the original Burnett equations and $`\alpha =\beta =1`$ the conventional Burnett equations. Let us now for simplicity denote nonlinear Burnett terms by dots. $$𝐏=p\mathrm{𝟏}2\mu 𝐒+\varpi _2\frac{\mu ^2}{p}\frac{D𝐯}{Dt}+\varpi _3\frac{\mu ^2}{\rho T}T+\mathrm{}$$ This gives $`𝐏`$ $`=`$ $`p\mathrm{𝟏}2\mu 𝐒+\varpi _2{\displaystyle \frac{\mu ^2}{p}}(1\alpha ){\displaystyle \frac{D𝐯}{Dt}}`$ $`+`$ $`{\displaystyle \frac{\mu ^2}{\rho T}}(\varpi _3\alpha \varpi _2)T\alpha \varpi _2{\displaystyle \frac{\mu ^2}{\rho ^2}}\rho +\mathrm{}`$ In order to study the linear stability of the resulting equations we now linearize the equations around a state at rest with constant temperature and density. $`\rho [1+\varpi _2{\displaystyle \frac{\mu ^2}{p\rho }}(1\alpha )({\displaystyle \frac{1}{2}}\mathrm{}+{\displaystyle \frac{1}{6}})]{\displaystyle \frac{𝐯}{t}}`$ $`=`$ $`{\displaystyle \frac{p}{\rho }}(1\alpha \varpi _2{\displaystyle \frac{\mu ^2}{\rho p}}{\displaystyle \frac{2}{3}}\mathrm{})\rho `$ $``$ $`{\displaystyle \frac{p}{T}}[1+{\displaystyle \frac{\mu ^2}{\rho p}}(\varpi _3\alpha \varpi _2){\displaystyle \frac{2}{3}}\mathrm{}]T+\mu [\mathrm{}𝐯+{\displaystyle \frac{1}{3}}(𝐯)].`$ All the undifferentiated quantities are here taken at the background state. For a plane wave with wave number $`k`$ the longitudinal component of the momentum equation gives $`{\displaystyle \frac{v_{}}{t}}`$ $`=`$ $`{\displaystyle \frac{1+\alpha \varpi _2\frac{\mu ^2}{\rho p}\frac{2}{3}k^2}{1+\varpi _2\frac{\mu ^2}{p\rho }(\alpha 1)\frac{2}{3}k^2}}{\displaystyle \frac{k_BT}{m\rho }}_{}\rho `$ $`{\displaystyle \frac{1+\frac{\mu ^2}{\rho p}(\alpha \varpi _2\varpi _3)\frac{2}{3}k^2}{1+\varpi _2\frac{\mu ^2}{p\rho }(\alpha 1)\frac{2}{3}k^2}}{\displaystyle \frac{k_B}{m}}_{}T`$ $`+{\displaystyle \frac{1}{1+\varpi _2\frac{\mu ^2}{p\rho }(\alpha 1)\frac{2}{3}k^2}}{\displaystyle \frac{4\mu }{3\rho }}\mathrm{}v_{}`$ Putting $`k=0`$ here we have the linearized Navier-Stokes equations. To avoid a singularity in the coefficients for a general $`k`$ we see that it is necessary that $`\alpha 1`$. As $`\alpha =0`$ corresponds to the original Burnett equations, we see that they have an unphysical singularity. - For a given $`k`$ we can interpret the equations as the linearized Navier-Stokes equations for an ideal gas but with different properties of the gas and the background state $`{\displaystyle \frac{\widehat{T}}{\widehat{m}}}`$ $`=`$ $`{\displaystyle \frac{T}{m}}{\displaystyle \frac{1+\alpha \varpi _2\frac{\mu ^2}{\rho p}\frac{2}{3}k^2}{1+\varpi _2\frac{\mu ^2}{p\rho }(\alpha 1)\frac{2}{3}k^2}},`$ $`{\displaystyle \frac{1}{\widehat{m}}}`$ $`=`$ $`{\displaystyle \frac{1}{m}}{\displaystyle \frac{1+\frac{\mu ^2}{\rho p}(\alpha \varpi _2\varpi _3)\frac{2}{3}k^2}{1+\varpi _2\frac{\mu ^2}{p\rho }(\alpha 1)\frac{2}{3}k^2}},`$ $`\widehat{\mu }`$ $`=`$ $`\mu {\displaystyle \frac{1}{[1+\varpi _2\frac{\mu ^2}{p\rho }(\alpha 1)\frac{2}{3}k^2]}}.`$ This interpretation is possible if the coefficients are positive. We see that this requires $$\alpha \varpi _2\varpi _30.$$ As $`0<\varpi _2<\varpi _3`$ for Maxwell molecules as well as hard spheres, $`\alpha `$ then has to be larger than $`1`$, which excludes the original Burnett equations as well as the conventional ones. Using $$𝐒=\frac{1}{2}\mathrm{}𝐯+\frac{1}{6}(𝐯)$$ we obtain for the heat current $`𝐪`$ $`=`$ $`\kappa T+\theta _2{\displaystyle \frac{\mu ^2}{\rho T}}(1\beta ){\displaystyle \frac{D(T)}{Dt}}`$ $`+{\displaystyle \frac{\mu ^2}{\rho }}[({\displaystyle \frac{\theta _4}{2}})\mathrm{}𝐯+({\displaystyle \frac{\theta _4}{6}}\beta {\displaystyle \frac{2\theta _2}{3}})(𝐯)]+\mathrm{}`$ The linearized energy equation is then $`{\displaystyle \frac{3k_B\rho }{2m}}(1+\theta _2{\displaystyle \frac{2}{3}}{\displaystyle \frac{\mu ^2}{\rho p}}(1\beta )\mathrm{}){\displaystyle \frac{T}{t}}`$ (13) $`=`$ $`p(1+{\displaystyle \frac{\mu ^2}{\rho p}}{\displaystyle \frac{2}{3}}(\theta _4\beta \theta _2)\mathrm{})(𝐯)+\kappa \mathrm{}T.`$ For a plane wave we obtain $`{\displaystyle \frac{T}{t}}={\displaystyle \frac{T}{mc_v}}{\displaystyle \frac{1+\frac{\mu ^2}{\rho p}\frac{2}{3}(\beta \theta _2\theta _4)k^2}{1+\theta _2\frac{2}{3}\frac{\mu ^2}{\rho p}(\beta 1)k^2}}(𝐯)`$ $`+{\displaystyle \frac{\kappa }{\rho c_v}}{\displaystyle \frac{1}{1+\theta _2\frac{2}{3}\frac{\mu ^2}{\rho p}(\beta 1)k^2}}\mathrm{}T.`$ Here we have introduced the specific heat $$c_v=\frac{3k_B}{2m}.$$ To avoid singularities it is necessary that $`\beta 1`$. This means that the conventional Burnett equations have no singularity in the energy equation, but the original Burnett equations do. We can interpret the energy equation as the energy equation for the Navier-Stokes equations of an ideal gas with Fourier’s expression for the heat current. $`{\displaystyle \frac{\widehat{T}}{\widehat{m}\widehat{c_v}}}`$ $`=`$ $`{\displaystyle \frac{T}{mc_v}}{\displaystyle \frac{1+\frac{\mu ^2}{\rho p}\frac{2}{3}(\beta \theta _2\theta _4)k^2}{1+\theta _2\frac{2}{3}\frac{\mu ^2}{\rho p}(\beta 1)k^2}},`$ $`{\displaystyle \frac{\widehat{\kappa }}{\widehat{c_v}}}`$ $`=`$ $`{\displaystyle \frac{\kappa }{c_v}}{\displaystyle \frac{1}{1+\theta _2\frac{2}{3}\frac{\mu ^2}{\rho p}(\beta 1)k^2}}.`$ In all we have then given new formal values to the five properties of the gas and its background state $`m,c_v,\mu ,\kappa ,T`$. $`\rho `$ and $`𝐯`$ are unchanged. This interpretation is possible as long as the coefficients are positive. As $`0<\theta _4<\theta _2`$ for Maxwell molecules as well as hard spheres this follows from $`\beta 1.`$ As the linearized equation of continuity is the usual one we conclude that for each value $`k`$ there is an ideal gas such that its linearized Navier-Stokes equations coincide with the full set of linearized hybrid Burnett equations for the same value of $`k`$. This means that the equations are linearly stable as long as $$\alpha \frac{\varpi _3}{\varpi _2},\beta 1.$$ (14) So we have found a two-parameter family of equations with linear stability. In this way of thinking, it is easy to understand that the conventional Burnett equations are linearly unstable. For large enough $`k`$ a temperature gradient gives rise to a pressure gradient in the opposite direction. So if a local region of the order of the mean free path or smaller is hotter than its surroundings, the pressure in this area will be lower and the gas will flow into it, increasing its temperature. But let us stress that this is for length scales of the order of the mean free path and smaller, so that the equations are not physically valid anyhow. As mentioned in the introduction, the equations nevertheless need to be well-behaved for such perturbations. For the transverse components we find $`\rho [1+\varpi _2{\displaystyle \frac{\mu ^2}{p\rho }}(\alpha 1){\displaystyle \frac{1}{2}}k^2]{\displaystyle \frac{𝐯_{}}{t}}`$ $`=`$ $`\mu \mathrm{}𝐯_{}`$ As long as $`\alpha 1`$ this is simply diffusion, with a diffusivity depending on $`k`$. ## IV Hybrid Burnett Equations Let us now make the following choice of the parameters $$\alpha =\frac{\varpi _3}{\varpi _2},\beta =1.$$ (15) The heat current then is the conventional one, $`𝐪_c`$. The resulting expression for the viscous pressure tensor is now denoted $`𝐏_h`$. $`𝐏_h=p\mathrm{𝟏}2\mu 𝐒+\varpi _1{\displaystyle \frac{\mu ^2}{p}}(𝐯)𝐒`$ $`+(\varpi _2\varpi _3){\displaystyle \frac{\mu ^2}{p}}{\displaystyle \frac{D𝐒}{Dt}}`$ $`2\varpi _2{\displaystyle \frac{\mu ^2}{p}}𝐒(𝐯)\varpi _3{\displaystyle \frac{\mu ^2}{\rho ^2}}\rho `$ $`+\varpi _3{\displaystyle \frac{\mu ^2}{\rho T}}{\displaystyle \frac{1}{\rho }}T\rho +{\displaystyle \frac{T}{\rho ^2}}\rho \rho {\displaystyle \frac{p}{\rho T}}(𝐯)^2`$ $`+\varpi _4{\displaystyle \frac{\mu ^2}{\rho pT}}pT+\varpi _5{\displaystyle \frac{\mu ^2}{\rho T^2}}TT`$ $`+\varpi _6{\displaystyle \frac{\mu ^2}{p}}𝐒𝐒.`$ Note that the troublesome third derivative of $`T`$ is absent. Further there is a change of sign in front of $`D𝐒/Dt`$ as compared to the original expression (II). Now we have our hybrid Burnett equations. $`{\displaystyle \frac{D\rho }{Dt}}+\rho 𝐯`$ $`=`$ $`0,`$ (16) $`\rho {\displaystyle \frac{D𝐯}{Dt}}`$ $`=`$ $`𝐏_h,`$ (17) $`\rho {\displaystyle \frac{3k_B}{2m}}{\displaystyle \frac{DT}{Dt}}`$ $`=`$ $`𝐏_h:𝐯𝐪_c.`$ (18) ## V Linear stability analysis Let us linearize around a uniform state at rest denoting now its temperature $`T_0`$ and density $`\rho _0`$. We write $$T=T_0(1+\stackrel{~}{T}),\rho =\rho _0(1+\stackrel{~}{\rho }),v=\sqrt{\frac{k_BT_0}{m}}\stackrel{~}{v}.$$ We introduce the dimensionless variables, where the unit of length is of the order of the mean free path $$x=x^{}\frac{\mu _0}{\rho _0}\sqrt{\frac{m}{k_BT_0}},t=t^{}\frac{\mu _0}{\rho _0}\frac{m}{k_BT_0}.$$ In the rest of this section stars and tildes are omitted and subscripts denote partial derivatives. We obtain in the longitudinal case $`\rho _t+v_x=0`$ $`\rho _x{\displaystyle \frac{2}{3}}\varpi _3\rho _{xxx}`$ $`+v_t{\displaystyle \frac{4}{3}}v_{xx}{\displaystyle \frac{2}{3}}(\varpi _3\varpi _2)v_{xxt}+T_x=0`$ $`{\displaystyle \frac{2}{3}}v_x+{\displaystyle \frac{4}{9}}(\theta _4\theta _2)v_{xxx}+T_tfT_{xx}=0`$ $`f=(2m\kappa )/(3k_B\mu )=5/(3Pr)`$, where $`Pr`$ is the Prandtl number. Looking for solutions $`\mathrm{exp}(ikx+\mathrm{\Lambda }t)`$, we find the determinant $`\mathrm{\Lambda }^3(1+{\displaystyle \frac{2\varpi _3}{3}}k^2)+\mathrm{\Lambda }^2[{\displaystyle \frac{2}{3}}(\varpi _3\varpi _2)fk^4+({\displaystyle \frac{4}{3}}+f)k^2]`$ $`+\mathrm{\Lambda }[({\displaystyle \frac{2\varpi _3}{3}}{\displaystyle \frac{4}{9}}(\theta _4\theta _2))k^4+{\displaystyle \frac{5}{3}}k^2]+(\varpi _3fk^6+fk^4).`$ Let us in particular consider the asymptotic case $`k\mathrm{}`$. We find one mode with $$\mathrm{\Lambda }_0\frac{(\varpi _3\varpi _2)f}{\varpi _3}k^2.$$ There are also two complex conjugate roots $`\mathrm{\Lambda }_\pm \pm i\sqrt{{\displaystyle \frac{3\varpi _3}{2(\varpi _3\varpi _2)}}}k`$ $`{\displaystyle \frac{9\varpi _3^2+2(\varpi _3\varpi _2)[3\varpi _3+2(\theta _2\theta _4)]}{6(\varpi _3\varpi _2)^2f}}.`$ For Maxwell molecules, see CC $$\theta _2=45/8,\theta _4=3,\varpi _2=2,\varpi _3=3.$$ For hard spheres, see Reinecke & Kremer $$\theta _2=5.826,\theta _4=2.416,\varpi _2=2.029,\varpi _3=2.415.$$ In both cases $`\varpi _3,\varpi _3\varpi _2,\theta _2\theta _4`$ are positive. This means that the mode corresponding to $`\mathrm{\Lambda }_0`$ is damped and nonpropagating. It is the entropy mode. The two modes $`\mathrm{\Lambda }_\pm `$ are damped. They are propagating sound waves. We plot the three roots $`\mathrm{\Lambda }`$ in the complex plane in Fig. 1. Here we use the value the Prandtl number $`2/3`$. This is the lowest approximation in terms of Sonine polynomial expansion for any interatomic potential and is experimentally found to be a good approximation, see CC . Prandtl numbers in the range between $`1/2`$ to $`1`$ give qualitatively the same plot as Fig. 1. We see in Fig. 1 that the real part of $`\mathrm{\Lambda }`$ is negativ. We have already in the preceding section found that the transverse mode is damped and nonpropagating. ## VI Discussion Let us pull out the acceleration from the viscous pressure tensor $`𝐏_h`$ using (7 ) $$𝐏_h=(\varpi _2\varpi _3)\frac{\mu ^2}{p}(\frac{D𝐯}{Dt})+\stackrel{~}{𝐏}_h.$$ Here $`\stackrel{~}{𝐏}_h=p\mathrm{𝟏}2\mu 𝐒+\varpi _1{\displaystyle \frac{\mu ^2}{p}}(𝐯)𝐒`$ $`(\varpi _2\varpi _3){\displaystyle \frac{\mu ^2}{p}}(𝐯)^2`$ $`2\varpi _2{\displaystyle \frac{\mu ^2}{p}}𝐒(𝐯)\varpi _3{\displaystyle \frac{\mu ^2}{\rho ^2}}\rho `$ $`+\varpi _3{\displaystyle \frac{\mu ^2}{\rho T}}{\displaystyle \frac{1}{\rho }}T\rho +{\displaystyle \frac{T}{\rho ^2}}\rho \rho {\displaystyle \frac{p}{\rho T}}(𝐯)^2`$ $`+\varpi _4{\displaystyle \frac{\mu ^2}{\rho pT}}pT+\varpi _5{\displaystyle \frac{\mu ^2}{\rho T^2}}TT`$ $`+\varpi _6{\displaystyle \frac{\mu ^2}{p}}𝐒𝐒.`$ The momentum equation can now be written $$\rho \frac{D𝐯}{Dt}[(\varpi _3\varpi _2)\frac{\mu ^2}{p}(\frac{D𝐯}{Dt})]=\stackrel{~}{𝐏}_h.$$ Let us write $`𝐚`$ for the acceleration and study the terms containing the acceleration $$\rho 𝐚[(\varpi _3\varpi _2)\frac{\mu ^2}{p}𝐚]$$ We multiply by $`𝐚`$ and integrate over the region of flow $`{\displaystyle \{𝐚[\rho 𝐚((\varpi _3\varpi _2)\frac{\mu ^2}{p}𝐚)]\}𝑑V}`$ $`=`$ $`{\displaystyle }[\rho a^2+(\varpi _3\varpi _2){\displaystyle \frac{\mu ^2}{p}}𝐚:𝐚]dV`$ $``$ $`{\displaystyle (\varpi _3\varpi _2)\frac{\mu ^2}{p}𝐚𝐚𝑑𝐒}`$ If there are no boundaries and the flow vanishes at infinity, the surface integral vanishes. We conclude that the operator acting on $`D𝐯/Dt`$ is then positive and the momentum equation can be solved for the acceleration. The situation when there are boundaries requires a closer examination. In order to see more clearly the structure of the hybrid Burnett equations, we replace the nonlinear terms in $`Kn^2`$ by dots. But first we split up in longitudinal and transverse parts according to $$\frac{1}{2}\mathrm{}+\frac{1}{6}=\frac{2}{3}+\frac{1}{2}(\mathrm{}).$$ The hybrid Burnett equations can then be written $`{\displaystyle \frac{D\rho }{Dt}}+\rho (𝐯)=0,`$ (19) $`\rho \{\mathrm{𝟏}(\varpi _3\varpi _2){\displaystyle \frac{\mu ^2}{p\rho }}[{\displaystyle \frac{1}{2}}(\mathrm{}\mathrm{𝟏})+{\displaystyle \frac{2}{3}}\}{\displaystyle \frac{D𝐯}{Dt}}`$ (20) $`=p+2(\mu 𝐒)+\varpi _3{\displaystyle \frac{\mu ^2}{\rho ^2}}{\displaystyle \frac{2}{3}}\mathrm{}\rho +\mathrm{},`$ $`{\displaystyle \frac{3k_B}{2m}}\rho {\displaystyle \frac{DT}{Dt}}=p(𝐯)+2\mu 𝐒𝐒+(\kappa T)`$ (21) $`(\theta _4\theta _2){\displaystyle \frac{\mu ^2}{\rho }}{\displaystyle \frac{2}{3}}\mathrm{}(𝐯)+\mathrm{}`$ Let us also write down the momentum equation of the conventional Burnett equations There is no need to write down the equations of continuity and energy as they are exactly the same as for the hybrid equations. $`\rho {\displaystyle \frac{D𝐯}{Dt}}=p+2(\mu 𝐒)(\varpi _3\varpi _2){\displaystyle \frac{\mu ^2}{\rho T}}{\displaystyle \frac{2}{3}}\mathrm{}T`$ (22) $`+\varpi _2{\displaystyle \frac{\mu ^2}{\rho ^2}}{\displaystyle \frac{2}{3}}\mathrm{}\rho +\mathrm{}`$ The difference is that the hybrid Burnett equations have a more complicated inertia term but no term $`\mathrm{}T`$ in the momentum equation. The coefficient in front of the term $`\mathrm{}\rho `$ is $`\varpi _3\varpi _2`$ in the conventional Burnett equations but $`\varpi _3`$ in the hybrid Burnett equations. Let us also note that the equations without the $`Kn^2`$ nonlinear terms are valid when $`Kn^2Ma^2`$ can be neglected and the relative temperature and density variations are assumed to be of the order of $`Ma`$. Let us take a closer look at the inertia term in the hybrid Burnett equations, see (17). We consider a Fourier component of the acceleration $`D𝐯/Dt`$ proportional to $`\mathrm{exp}(i𝐤𝐫)`$. Then $`\rho \{[\mathrm{𝟏}(\varpi _3\varpi _2){\displaystyle \frac{\mu ^2}{p\rho }}[{\displaystyle \frac{1}{2}}(\mathrm{}\mathrm{𝟏})+{\displaystyle \frac{2}{3}}]\}{\displaystyle \frac{D𝐯}{Dt}}`$ $`=`$ $`\rho \{[1+{\displaystyle \frac{1}{2}}(\varpi _3\varpi _2){\displaystyle \frac{\mu ^2k^2}{p\rho }}](\mathrm{𝟏}{\displaystyle \frac{\mathrm{𝐤𝐤}}{k^2}})`$ $`+[1+{\displaystyle \frac{2}{3}}(\varpi _3\varpi _2){\displaystyle \frac{\mu ^2k^2}{p\rho }}]{\displaystyle \frac{\mathrm{𝐤𝐤}}{k^2}}\}{\displaystyle \frac{D𝐯}{Dt}}.`$ We have effectively a transverse inertia in the first term and a longitudinal inertia in the second term. As $`\varpi _3\varpi _2>0`$, both of them are positive. We noted earlier that the sign of the $`D𝐯/Dt`$ term of the hybrid Burnett expression for $`𝐏`$ is the opposite of that of the original Burnett equations. As a result the original Burnett equations have a singularity for certain wave numbers where the inertia vanishes. Let us also consider the low $`Ma`$ stationary case. The hybrid Burnett equations are then $`(\rho 𝐯)`$ $`=`$ $`0,`$ $`\rho (𝐯)v=p`$ $`+`$ $`2(\mu 𝐒)+\varpi _3{\displaystyle \frac{\mu ^2}{\rho ^2}}{\displaystyle \frac{2}{3}}\mathrm{}\rho ,`$ $`{\displaystyle \frac{3k_B}{2m}}\rho (𝐯T)`$ $`=`$ $`p(𝐯)+2\mu 𝐒:𝐒`$ $`+(\kappa T)`$ $``$ $`(\theta _4\theta _2){\displaystyle \frac{\mu ^2}{\rho }}{\displaystyle \frac{2}{3}}\mathrm{}(𝐯).`$ The conventional Burnett momentum equation is $`\rho (𝐯)𝐯=p+2(\mu 𝐒)`$ (23) $`(\varpi _3\varpi _2){\displaystyle \frac{\mu ^2}{\rho T}}{\displaystyle \frac{2}{3}}\mathrm{}T+\varpi _2{\displaystyle \frac{\mu ^2}{\rho ^2}}{\displaystyle \frac{2}{3}}\mathrm{}\rho ,`$ The only difference between the hybrid Burnett equations and the conventional Burnett equations is the change of coefficients $$\varpi _2\varpi _3,\varpi _3\varpi _3.$$ (24) This also means that the third derivative of the temperature is absent in the hybrid Burnett equations. ###### Acknowledgements. During the years I have had many stimulating discussions with Prof. Y. Sone on the relation between the asymptotic method and the Chapman-Enskog expansion.
warning/0506/math0506600.html	ar5iv	text	# Associativity as Commutativity ## 1 Introduction Associativity is a kind of commutativity. To see why, conceive of $`(a(cb))`$ as $`(a\underset{¯}{})(\underset{¯}{}b)`$ applied to $`c`$. We have \| $`((a\underset{¯}{})(\underset{¯}{}b))(c)`$ \| $`=(a\underset{¯}{})((\underset{¯}{}b)(c))`$ \| \| --- \| --- \| \| \| $`=(a\underset{¯}{})((cb))`$ \| \| \| $`=(a(cb)).`$ \| Then associativity goes from $`(a\underset{¯}{})(\underset{¯}{}b)`$ to $`(\underset{¯}{}b)(a\underset{¯}{})`$, which when applied to $`c`$ yields $`((ac)b)`$. The purpose of this paper is to exploit this idea to show that monoidal categories may be conceived as a kind of symmetric strictly monoidal categories, where associativity arrows are identities. As a matter of fact, this analogy holds also with braided strictly monoidal categories. More precisely, we show that coherence conditions for monoidal categories concerning associativity are analogous to coherence conditions concerning commutativity (i.e. symmetry or braiding) for symmetric (see , or ) or braided (see and , second edition, Chapter IX) strictly monoidal categories. In particular, Mac Lane’s pentagonal coherence condition for associativity (see Section 4 below) is decomposed into conditions concerning commutativity, among which we have a condition analogous to naturality and degenerate cases of Mac Lane’s hexagonal condition for commutativity. (The hexagon becomes a triangle, because associativity arrows are identities, or a two-sided figure.) This decomposition is analogous to the derivation of the Yang-Baxter equation from Mac Lane’s hexagon and the naturality of commutativity (see Section 5 below). To achieve that, we replace the algebra freely generated with one binary operation (denoted by $``$) by an isomorphic algebra generated with a family of partial operations we call insertion (denoted by $`_n`$); insertion is analogous to the composition $``$ at the beginning of this text, or to functional application. (This procedure is like Achilles’ introduction of $`\alpha `$ in .) The latter algebra is more complicated, and is not free any more, but it enables us to present associativity arrows as commutativity arrows. In the next section we state precisely these matters concerning insertion. After that we introduce a category $`𝚪`$, which in the remainder of the paper is shown isomorphic to a free monoidal category without unit. In $`𝚪`$ the associativity arrows appear as a kind of commutativity arrows, and coherence conditions for $`𝚪`$ take the form of an inductive definition of these commutativity arrows. Our decomposition of Mac Lane’s pentagon is in the last section. We work with categories where associativity is an isomorphism, because this is the standard approach, but our treatment is easily transferred to categories where associativity arrows are not necessarily isomorphisms, which in (Section 4.2) are called semiassociative categories (see also ). As monoidal categories, semiassociative categories are coherent in Mac Lane’s “all diagrams commute” sense. With semiassociative categories, the commutativity corresponding to associativity only ceases to be an isomorphism, and all the rest remains as in the text that follows. Among the coherence conditions for the main kinds of categories with structure, Mac Lane’s pentagon seems to be more mysterious than the others. Our decomposition of the pentagon goes towards dispelling the mystery. The pentagon is reduced to an inductive definition of a kind of commutativity. If the associativity arrows are isomorphisms, then the pentagon yields the definition of an associativity arrow complex in one of its indices in terms of associativity arrows simpler in that index, but more complex in the other two indices. There is no reduction of complexity in all the indices, and no real inductive definition. Our approach, which works also in the absence of isomorphism, as we said above, gives a real inductive definition. ## 2 Insertion Let $`_1`$ be the set of words (finite sequences of symbols) in the alphabet $`\{\text{},,(,)\}`$ defined inductively by \| $`\text{}_1`$, \| \| --- \| \| if $`X,Y_1`$, then $`(XY)_1`$. \| Let $`_2`$ be $`_1\{\text{}\}`$. The elements of $`_2`$ may be identified with finite planar binary trees with more than one node, while is the trivial one-node tree. In this section, we use $`X`$, $`Y`$ and $`Z`$ for the members of $`_1`$. (Starting from the end of the section, we change this notation to $`A`$, $`B`$, $`C,\mathrm{}`$) We omit the outermost pair of parentheses of the members of $`_1`$, taking them for granted. We make the same omission in other analogous situations later on. Let $`𝐍^+`$ be the set of natural numbers greater than $`0`$, and let $`^{}`$ be the set of words in the alphabet $`\{\mathrm{𝟏},\mathrm{𝟐}\}\{_n\|n𝐍^+\}`$ defined inductively by the following clauses that involve also an inductive definition of a map $`\|\|`$ from $`^{}`$ to $`𝐍^+`$: \| \| $`\mathrm{𝟏}^{}`$ and $`\|\mathrm{𝟏}\|=1`$, \| \| --- \| --- \| \| \| $`\mathrm{𝟐}^{}`$ and $`\|\mathrm{𝟐}\|=2`$, \| \| if $`A,B^{}`$ and $`1n\|A\|`$, then \| \| \| $`(A_nB)^{}`$ and $`\|(A_nB)\|=\|A\|+\|B\|1`$. \| Let $`^{\prime \prime }`$ be defined as $`^{}`$ save that we omit the first clause above involving $`\mathrm{𝟏}`$, and we replace $`^{}`$ by $`^{\prime \prime }`$ in the two remaining clauses. For the members of $`^{}`$ we use $`A`$, $`B`$, $`C,\mathrm{}`$, sometimes with indices. As we did for $`_1`$, we omit the outermost pair of parentheses of the members of $`^{}`$. We define the equational calculus $`^{\prime \prime }`$ in $`^{\prime \prime }`$ (i.e. a calculus whose theorems are equations between members of $`^{\prime \prime }`$) by assuming reflexivity, symmetry and transitivity of equality, the rule that if $`A=B`$ and $`C=D`$, then $`A_nC=B_nD`$, provided $`A_nC`$ and $`B_nD`$ are defined, and the two axioms \| (assoc 1) \| $`(A_nB)_mC=A_n(B_{mn+1}C)`$ \| if $`nm<n+\|B\|`$, \| \| --- \| --- \| --- \| \| (assoc 2) \| $`(A_nB)_mC=(A_{m\|B\|+1}C)_nB`$ \| if $`n+\|B\|m`$. \| Note that the condition $`nm<n+\|B\|`$ in (assoc 1) follows from the legitimacy of $`B_{mn+1}C`$. Note also that in both (assoc 1) and (assoc 2) we have $`nm`$. The equation (assoc 2) could be replaced by \| $`(A_nB)_mC=(A_mC)_{n+\|C\|1}B`$ \| if $`m<n`$. \| \| --- \| --- \| (The equations (assoc 1) and (assoc 2) are analogous to the two associativity equations for the cut operation one finds in multicategories; see and , Section 3. Analogous equations are also found in the definition of operad; see , Section 1.) The equational calculus $`^{}`$ in $`^{}`$ is defined as $`^{\prime \prime }`$ with the additional axiom \| (unit) \| $`\mathrm{𝟏}_1A=A_n\mathrm{𝟏}=A`$ \| \| --- \| --- \| (whose analogue one also finds in multicategories). Our purpose now is to interpret $`^{}`$ in $`_1`$. This will make clear the meaning of the axioms of $`^{}`$. For $`X`$ in $`_1`$, let $`\|X\|`$ be the number of occurrences of in $`X`$. We define in $`_1`$ the partial operation of insertion $`_n`$ by the following inductive clauses: \| $`\text{}_1Z=Z`$, \| \| --- \| \| $`(XY)_nZ=\{\begin{array}{cc}(X_nZ)Y\hfill & \text{if }1n\|X\|\hfill \\ X(Y_{n\|X\|}Z)\hfill & \text{if }\|X\|<n\|X\|+\|Y\|.\hfill \end{array}`$ \| We define insertion in $`_2`$ by replacing the clause $`\text{}_1Z=Z`$ above by the clauses $$\begin{array}{c}(\text{}\text{})_1Z=Z\text{},\hfill \\ (\text{}\text{})_2Z=\text{}Z.\hfill \end{array}$$ Insertion gets its name from the fact that $`X_nZ`$ is obtained by inserting $`Z`$ at the place of the $`n`$-th occurrence of in $`X`$, starting from the left; namely, the $`n`$-th leaf of the tree corresponding to $`X`$ becomes the root of the tree corresponding to $`Z`$, and the resulting tree corresponds to $`X_nZ`$. Insertion is called grafting in , and particular instances of insertion, which one finds in the source and target of the arrows $`\gamma _{A,B}^{}`$ in Section 3 below, are called under and over in (Section 1.5). We interpret $`^{\prime \prime }`$ in $`_2`$, i.e., we define a function $`v`$ from $`^{\prime \prime }`$ to $`_2`$, in the following manner: $$\begin{array}{c}v(\mathrm{𝟐})=\text{}\text{},\hfill \\ v(A_nB)=v(A)_nv(B).\hfill \end{array}$$ For this definition to be correct, we must check that $`\|A\|=\|v(A)\|`$, which is easily done by induction on the length of $`\|A\|`$. We prove first the following by an easy induction on the length of derivation. Soundness. If $`A=B`$ in $`^{\prime \prime }`$, then $`v(A)=v(B)`$. Our purpose is to prove also the converse: Completeness. If $`v(A)=v(B)`$, then $`A=B`$ in $`^{\prime \prime }`$. For every $`A`$ in $`^{\prime \prime }`$ we define the natural number $`c(A)`$ inductively as follows: $$\begin{array}{c}c(\mathrm{𝟐})=2,\hfill \\ c(B_nC)=c(B)(c(C)+1).\hfill \end{array}$$ Let $`s(A)`$ be the sum of the indices $`n`$ of all the occurrences of $`_n`$ in $`A`$, and let $`d(A)=c(A)+s(A)`$. Then we can easily check that if $`A=B`$ is an instance of (assoc 1) or (assoc 2), then $`d(A)>d(B)`$. Let a member of $`^{\prime \prime }`$ be called normal when it has no part of the form of the left-hand side of (assoc 1) or (assoc 2), i.e. no part of the form $`(A_nB)_mC`$ for $`nm`$. It can be shown that a normal member of $`^{\prime \prime }`$ is of one of the following forms: $$(\mathrm{𝟐}_2A_2)_1A_1,\mathbf{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}_2A_2,\mathbf{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}_1A_1,\mathbf{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2},$$ for $`A_1`$ and $`A_2`$ normal. These are the four normal types. Then it is easy to show by applying (assoc 1) and (assoc 2) from left to right that for every $`A`$ in $`^{\prime \prime }`$ there is a normal $`A^{}`$ such that $`A=A^{}`$ in $`^{\prime \prime }`$. We can also show the following. Auxiliary Lemma. If $`A`$ and $`B`$ are normal and $`v(A)=v(B)`$, then $`A`$ and $`B`$ coincide. Proof. If $`v(A)=v(B)`$, then $`A`$ and $`B`$ must be of the same normal type (otherwise, clearly, $`v(A)v(B)`$). If $`A`$ is $`(\mathrm{𝟐}_2A_2)_1A_1`$ and $`B`$ is $`(\mathrm{𝟐}_2B_2)_1B_1`$, then we conclude that $`v(A_1)=v(B_1)`$ and $`v(A_2)=v(B_2)`$, and we reason by induction. We reason analogously for the second and third normal type, and the normal type $`\mathrm{𝟐}`$ provides the basis of the induction. $``$ To prove Completeness, suppose $`v(A)=v(B)`$. Let $`A=A^{}`$ and $`B=B^{}`$ in $`^{\prime \prime }`$ for $`A^{}`$ and $`B^{}`$ normal. Then by Soundness we have $`v(A^{})=v(A)=v(B)=v(B^{})`$, and so, by the Auxiliary Lemma, $`A^{}`$ and $`B^{}`$ coincide. It follows that $`A=B`$ in $`^{\prime \prime }`$, which proves Completeness. We can now also show that if $`A=A^{}`$ and $`A=A^{\prime \prime }`$ in $`^{\prime \prime }`$ for $`A^{}`$ and $`A^{\prime \prime }`$ normal, then $`A^{}`$ and $`A^{\prime \prime }`$ coincide. This follows from Soundness and the Auxiliary Lemma. (One could show this uniqueness of normal form directly in $`^{\prime \prime }`$, without proceeding via $`v`$ and $`_2`$, by relying on confluence techniques, as in the lambda calculus or term-rewriting systems. In such a proof, diagrams analogous to Mac Lane’s pentagon and the Yang-Baxter equation would arise.) We interpret $`^{}`$ in $`_1`$ by extending the definition of $`v`$ from $`^{\prime \prime }`$ to $`_2`$ with the clause $`v(\mathrm{𝟏})=\text{}`$. Then we can prove Soundness and Completeness with $`^{\prime \prime }`$ replaced by $`^{}`$. In reducing a member of $`^{}`$ to a normal member of $`^{\prime \prime }`$ or to $`\mathrm{𝟏}`$ we get rid first of all superfluous occurrences of $`\mathrm{𝟏}`$, by relying on the equations (unit). Otherwise, the proof proceeds as before. We can factorize $`^{}`$ through the smallest equivalence relation such that the equations of $`^{}`$ are satisfied, and obtain a set of equivalence classes isomorphic to $`_1`$. For the equivalence classes $`[A]`$ and $`[B]`$ we define $``$ by $$[A][B]=_{\text{df}}[(\mathrm{𝟐}_2B)_1A],$$ and the isomorphism $`i`$ from $`_1`$ to $`^{}`$ is defined by $$\begin{array}{c}i(\text{})=[\mathrm{𝟏}],\hfill \\ i(XY)=i(X)i(Y).\hfill \end{array}$$ The inverse $`i^1`$ of $`i`$ is defined by $`i^1([A])=v(A)`$. (To verify that $`i`$ and $`i^1`$ are inverse to each other, we rely on the fact that every $`[A]`$ is equal to $`[A^{}]`$ for $`A^{}`$ being normal or $`\mathrm{𝟏}`$.) We designate the equivalence class $`[A]`$ by $`A`$, and to designate the elements of $`_1`$ we can then use the notation introduced for $`^{}`$. This means that we can write $`A`$, $`B`$, $`C,\mathrm{}`$ for $`X`$, $`Y`$, $`Z,\mathrm{}`$, we can write $`\mathrm{𝟐}`$ for $`\text{}\text{}`$, and we can write $`_n`$ for $`_n`$. We have for $`_1`$ the equation $$AB=(\mathrm{𝟐}_2B)_1A.$$ We will write however instead of 1, and reserve 1 with a subscript for the name of an arrow. ## 3 The category $`𝚪`$ The objects of the category $`𝚪`$ are the elements of $`_1`$. To define the arrows of $`𝚪`$, we define first inductively the arrow terms of $`𝚪`$ in the following way: $$\begin{array}{c}\mathrm{𝟏}_A:AA,\hfill \\ \gamma _{A,B}^{}:A_{\|A\|}BB_1A,\hfill \\ \gamma _{A,B}^{}:B_1AA_{\|A\|}B\hfill \end{array}$$ are arrow terms of $`𝚪`$ for all objects $`A`$ and $`B`$; if $`f:AB`$ and $`g:CD`$ are arrow terms of $`𝚪`$, then $`gf:AD`$ is an arrow term of $`𝚪`$, provided $`B`$ is $`C`$, and $`f_ng:A_nCB_nD`$ is an arrow term of $`𝚪`$, provided $`1n\|A\|`$ and $`1n\|B\|`$. Note that for all arrow terms $`f:AB`$ of $`𝚪`$ we have $`\|A\|=\|B\|`$; we write $`\|f\|`$ for $`\|A\|`$, which is equal to $`\|B\|`$. The arrows of $`𝚪`$ are equivalence classes of arrow terms of $`𝚪`$ (cf. , Section 2.3) such that the following equations are satisfied: \| \| (cat 1) \| $`\mathrm{𝟏}_Bf=f\mathbf{\hspace{0.17em}1}_A=f`$, for $`f:AB`$, \| \| --- \| --- \| --- \| \| \| (cat 2) \| $`(hg)f=h(gf)`$, \| \| \| (bif 1) \| $`\mathrm{𝟏}_A_n\mathrm{𝟏}_B=\mathrm{𝟏}_{A_nB}`$, \| \| \| (bif 2) \| $`(f_2f_1)_n(g_2g_1)=(f_2_ng_2)(f_1_ng_1)`$, \| \| for $`1n\|f\|`$ and $`nm\|f\|+\|g\|1`$ \| \| \| (assoc 1$``$) \| $`(f_ng)_mh=f_n(g_{mn+1}h)`$ \| if $`nm<n+\|g\|`$, \| \| \| (assoc 2$``$) \| $`(f_ng)_mh=(f_{m\|g\|+1}h)_ng`$ \| if $`n+\|g\|m`$, \| \| \| (unit $``$) \| $`\mathrm{𝟏}_{\text{}}_1f=f_n\mathrm{𝟏}_{\text{}}=f`$, \| \| \| ($`\gamma `$ nat) \| $`\gamma _{B,D}^{}(f_{\|A\|}g)=(g_1f)\gamma _{A,C}^{}`$, \| \| \| ($`\gamma \gamma `$) \| $`\gamma _{A,B}^{}\gamma _{A,B}^{}=\mathrm{𝟏}_{A_{\|A\|}B}`$, $`\gamma _{A,B}^{}\gamma _{A,B}^{}=\mathrm{𝟏}_{B_1A}`$, \| \| \| ($`\gamma `$1) \| $`\gamma _{\text{},A}^{}=\gamma _{A,\text{}}^{}=\mathrm{𝟏}_A`$, \| \| \| (hex 1) \| $`\gamma _{A_{\|A\|}B,C}^{}=(\gamma _{A,C}^{}_{\|A\|}\mathrm{𝟏}_B)(\mathrm{𝟏}_A_{\|A\|}\gamma _{B,C}^{})`$, \| \| \| (hex 1a) \| $`\gamma _{A_nB,C}^{}=\gamma _{A,C}^{}_n\mathrm{𝟏}_B`$ \| if $`1n<\|A\|`$, \| \| \| (hex 2) \| $`\gamma _{C,A_1B}^{}=(\mathrm{𝟏}_A_1\gamma _{C,B}^{})(\gamma _{C,A}^{}_{\|C\|}\mathrm{𝟏}_B)`$, \| \| \| (hex 2a) \| $`\gamma _{C,A_nB}^{}=\gamma _{C,A}^{}_{n+\|C\|1}\mathrm{𝟏}_B`$ \| if $`1<n\|A\|`$. \| We also assume besides reflexivity, symmetry and transitivity of equality that if $`f=g`$ and $`h=j`$, then for $`\alpha `$ being $``$ or $`_n`$ we have $`f\alpha h=g\alpha j`$, provided $`f\alpha h`$ and $`g\alpha j`$ are defined. The equations (cat 1) and (cat 2) make of $`𝚪`$ a category. The equations (bif 1) and (bif 2) are analogous to bifunctorial equations. The equations (assoc 1$``$), (assoc 2$``$) and (unit $``$) are analogous to naturality equations. In (assoc 1$``$) the associativity arrows with respect to $`_n`$ are not written down because they are identity arrows, in virtue of the equation (assoc 1) on objects. Analogous remarks hold for (assoc 2$``$) and (unit $``$). The equation ($`\gamma `$ nat) is analogous to a naturality equation, and the equations ($`\gamma \gamma `$) say that $`\gamma _{A,B}^{}`$ is an isomorphism, with inverse $`\gamma _{A,B}^{}`$. The equation ($`\gamma \mathrm{𝟏}`$) is auxiliary, and would not be needed if we had assumed $`\gamma _{A,B}^{}`$ and $`\gamma _{A,B}^{}`$ only for $`A`$ and $`B`$ different from . The equations (hex 1) and (hex 2) are analogous to Mac Lane’s hexagonal equation of symmetric monoidal categories (see , , Section VII.7, or , Section 5.1). Here the associativity arrows with respect to $`_n`$ are identity arrows, in virtue of the equation (assoc 1) on objects (and so instead of hexagons we have triangles; cf. the equation ($`c`$ hex 1) in Section 5 below). Finally, the equations (hex 1a) and (hex 2a), together with ($`\gamma \mathrm{𝟏}`$), (hex 1) and (hex 2), enable us to define inductively $`\gamma _{A,B}^{}`$ for all $`A`$ and $`B`$ in terms of the identity arrows $`\mathrm{𝟏}_A`$, the arrows $$\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}:\text{}(\text{}\text{})(\text{}\text{})\text{}$$ and the operations on arrows $``$ and $`_n`$. Relying on ($`\gamma \gamma `$), we can proceed analogously for $`\gamma _{A,B}^{}`$ by using instead of $`\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}`$ the arrows $$\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}:(\text{}\text{})\text{}\text{}(\text{}\text{}).$$ The equations (hex 1a) and (hex 2a) are also analogous to Mac Lane’s hexagon mentioned above (due to the presence of (assoc 2) too, the collapse is however not any more into a triangle, but into a two-sided figure). ## 4 The category Â The category Â has the same objects as $`𝚪`$; namely, the elements of $`_1`$. To define the arrows of Â, we define first inductively the arrow terms of Â in the following way: $$\begin{array}{c}\mathrm{𝟏}_A:AA,\hfill \\ b_{A,B,C}^{}:A(BC)(AB)C,\hfill \\ b_{A,B,C}^{}:(AB)CA(BC)\hfill \end{array}$$ are arrow terms of Â for all objects $`A`$, $`B`$ and $`C`$; if $`f:AB`$ and $`g:CD`$ are arrow terms of Â, then $`gf:AD`$ is an arrow term of Â, provided $`B`$ is $`C`$, and $`fg:ACBD`$ is an arrow term of Â. The arrows of Â are equivalence classes of arrow terms of Â such that the following equations are satisfied: (cat 1), (cat 2), (bif 1) and (bif 2) with $`_n`$ replaced by $``$, and moreover \| \| ($`b`$ nat) \| $`b_{B,D,F}^{}(f(gh))=((fg)h)b_{A,C,E}^{}`$, \| \| --- \| --- \| --- \| \| \| ($`bb`$) \| $`b_{A,B,C}^{}b_{A,B,C}^{}=\mathrm{𝟏}_{A(BC)}`$, $`b_{A,B,C}^{}b_{A,B,C}^{}=\mathrm{𝟏}_{(AB)C}`$, \| \| \| ($`b`$5) \| $`b_{AB,C,D}^{}b_{A,B,CD}^{}=(b_{A,B,C}^{}\mathrm{𝟏}_D)b_{A,BC,D}^{}(\mathrm{𝟏}_Ab_{B,C,D}^{})`$. \| We also assume besides reflexivity, symmetry and transitivity of equality that if $`f=g`$ and $`h=j`$, then $`fh=gj`$, provided $`fh`$ and $`gj`$ are defined, and $`fh=gj`$. In Â we have that $``$ is a bifunctor, $`b^{}`$ is a natural isomorphism in all its indices, and ($`b`$5) is Mac Lane’s pentagonal equation of , where it is proved that the category Â is a preorder. Namely, for all arrows $`f,g:AB`$ of Â we have that $`f=g`$ (see also , Section VII.2, or , Section 4.3). The category Â is the free monoidal category without unit, i.e. free associative category in the terminology of (Section 4.3), generated by a single object, this object being conceived as a trivial discrete category. ## 5 The isomorphism of $`𝚪`$ and Â We are going to prove that the categories $`𝚪`$ and Â are isomorphic. We define first what is missing of the structure of Â in $`𝚪`$ in the following manner: \| \| $`b_{A,B,C}^{}`$ \| $`=_{\text{df}}`$ \| $`((\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_3\mathrm{𝟏}_C)_2\mathrm{𝟏}_B)_1\mathrm{𝟏}_A`$, \| \| --- \| --- \| --- \| --- \| \| \| $`b_{A,B,C}^{}`$ \| $`=_{\text{df}}`$ \| $`((\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_3\mathrm{𝟏}_C)_2\mathrm{𝟏}_B)_1\mathrm{𝟏}_A`$, \| \| \| \| $`fg`$ $`=_{\text{df}}`$ \| $`(\mathrm{𝟏}_\mathrm{𝟐}_2g)_1f`$. \| It can then be checked by induction on the length of derivation that the equations of Â are satisfied in $`𝚪`$. We have of course the equations (cat 1) and (cat 2), while the equations (bif 1) and (bif 2) with $`_n`$ replaced by $``$ are easy consequences of (bif 1) and (bif 2). To derive ($`b`$ nat), we have that with the help of (assoc 1$``$) and (bif 1) the left-hand side is equal to $$(((\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_3\mathrm{𝟏}_F)_2\mathrm{𝟏}_D)_1\mathrm{𝟏}_B)(((\mathrm{𝟏}_{\mathrm{𝟐}_2\mathrm{𝟐}}_3h)_2g)_1f),$$ while with the help of (assoc 1$``$), (assoc 2$``$) and (bif 1) the right-hand side is equal to $$(((\mathrm{𝟏}_{\mathrm{𝟐}_1\mathrm{𝟐}}_3h)_2g)_1f)(((\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_3\mathrm{𝟏}_E)_2\mathrm{𝟏}_C)_1\mathrm{𝟏}_A).$$ Then it is enough to apply (bif 2) and (cat 1). It is trivial to derive ($`bb`$) with the help of (bif 2), ($`\gamma \gamma `$) and (bif 1). We derive finally the pentagonal equation ($`b`$5). With the help of (bif 1), (assoc 1$``$) and (assoc 2$``$) we derive that each of $$b_{A,B,CD}^{},b_{AB,C,D}^{},\mathbf{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}_Ab_{B,C,D}^{},b_{A,BC,D}^{},b_{A,B,C}^{}\mathrm{𝟏}_D$$ is equal to $`(((f_4\mathrm{𝟏}_D)_3\mathrm{𝟏}_C)_2\mathrm{𝟏}_B)_1\mathrm{𝟏}_A`$ for $`f`$ being respectively $$\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_3\mathrm{𝟏}_\mathrm{𝟐},\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_1\mathrm{𝟏}_\mathrm{𝟐},\mathbf{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}_\mathrm{𝟐}_2\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{},\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_2\mathrm{𝟏}_\mathrm{𝟐},\mathbf{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}_\mathrm{𝟐}_1\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}.$$ Then, by relying on (bif 2), it is enough to derive the following: \| $`(\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_1\mathrm{𝟏}_\mathrm{𝟐})(\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_3\mathrm{𝟏}_\mathrm{𝟐})=\gamma _{\mathrm{𝟐}_1\mathrm{𝟐},\mathrm{𝟐}}^{}(\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_3\mathrm{𝟏}_\mathrm{𝟐})`$, by (hex 1a), \| \| --- \| \| \| = \| $`(\mathrm{𝟏}_\mathrm{𝟐}_1\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{})\gamma _{\mathrm{𝟐}_2\mathrm{𝟐},\mathrm{𝟐}}^{}`$, by ($`\gamma `$ nat), \| \| \| = \| $`(\mathrm{𝟏}_\mathrm{𝟐}_1\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{})(\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_2\mathrm{𝟏}_\mathrm{𝟐})(\mathrm{𝟏}_\mathrm{𝟐}_2\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{})`$, by (hex 1). \| Diagrammatically, we have So the pentagon is decomposed into a triangle (a degenerate hexagon, corresponding to (hex 1)), a square (analogous to a naturality square, corresponding to ($`\gamma `$ nat)) and a two-sided diagram (corresponding to (hex 1a)). If $`c_{A,B}:ABBA`$ is the commutativity arrow of symmetric monoidal categories, for which in strict categories of this kind, where associativity arrows are identities, we have the equations \| ($`c`$ nat) \| $`c_{B,D}(fg)=(gf)c_{A,B}`$, \| \| --- \| --- \| \| ($`c`$ hex 1) \| $`c_{AB,C}=(c_{A,C}\mathrm{𝟏}_B)(\mathrm{𝟏}_Ac_{B,C})`$, \| then we derive the Yang-Baxter equation in the following way: \| $`(c_{B,C}\mathrm{𝟏}_A)(\mathrm{𝟏}_Bc_{A,C})(c_{A,B}\mathrm{𝟏}_C)=c_{BA,C}(c_{A,B}\mathrm{𝟏}_C)`$, by ($`c`$ hex 1), \| \| --- \| \| $`=(\mathrm{𝟏}_Cc_{A,B})c_{AB,C}`$, by ($`c`$ nat), \| \| $`=(\mathrm{𝟏}_Cc_{A,B})(c_{A,C}\mathrm{𝟏}_B)(\mathrm{𝟏}_Ac_{B,C})`$, by ($`c`$ hex 1). \| This derivation is analogous to our derivation of ($`b`$5) above, where however the arrow corresponding to $`\mathrm{𝟏}_Bc_{A,C}`$ is identity, in virtue of the equation (assoc 2) on objects. Alternatively, we derive ($`b`$5) by using the following: \| $`(\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_1\mathrm{𝟏}_\mathrm{𝟐})(\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_3\mathrm{𝟏}_\mathrm{𝟐})=(\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_1\mathrm{𝟏}_\mathrm{𝟐})\gamma _{\mathrm{𝟐},\mathrm{𝟐}_2\mathrm{𝟐}}^{}`$, by (hex 2a), \| \| --- \| \| \| = \| $`\gamma _{\mathrm{𝟐},\mathrm{𝟐}_1\mathrm{𝟐}}^{}(\mathrm{𝟏}_\mathrm{𝟐}_2\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{})`$, by ($`\gamma `$ nat), \| \| \| = \| $`(\mathrm{𝟏}_\mathrm{𝟐}_1\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{})(\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}_2\mathrm{𝟏}_\mathrm{𝟐})(\mathrm{𝟏}_\mathrm{𝟐}_2\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{})`$, by (hex 2). \| Diagrammatically, we have This is an alternative decomposition of the pentagon into a triangle, a square and a two-sided diagram. Hence we have in $`𝚪`$ all the equations of Â. To define what is missing of the structure of $`𝚪`$ in Â, we have first the following inductive definition of $`_n`$ on arrows: \| if $`A^{}(B^{}C^{})=(A(BC))_nD`$, \| \| --- \| \| \| $`b_{A,B,C}^{}_n\mathrm{𝟏}_D=b_{A^{},B^{},C^{}}^{}`$, $`b_{A,B,C}^{}_n\mathrm{𝟏}_D=b_{A^{},B^{},C^{}}^{}`$, \| \| \| $`(gf)_n\mathrm{𝟏}_D=(g_n\mathrm{𝟏}_D)(f_n\mathrm{𝟏}_D)`$, \| \| \| $`(fg)_n\mathrm{𝟏}_D=\{\begin{array}{cc}(f_n\mathrm{𝟏}_D)g\hfill & \text{if }1n\|f\|\hfill \\ f(g_{n\|f\|}\mathrm{𝟏}_D)\hfill & \text{if }\|f\|<n\|f\|+\|g\|,\hfill \end{array}`$ \| \| \| $`\mathrm{𝟏}_{\text{}}_1f=f`$, \| \| \| $`\mathrm{𝟏}_{AB}_nf=\{\begin{array}{cc}(\mathrm{𝟏}_A_nf)\mathrm{𝟏}_B\hfill & \text{if }1n\|A\|\hfill \\ \mathrm{𝟏}_A(\mathrm{𝟏}_B_{n\|A\|}f)\hfill & \text{if }\|A\|<n\|A\|+\|B\|,\hfill \end{array}`$ \| \| \| $`f_ng=(\mathrm{𝟏}_B_ng)(f_n\mathrm{𝟏}_C)`$. \| We define $`\gamma _{A,B}^{}`$ and $`\gamma _{A,B}^{}`$ by stipulating $$\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}=_{\text{df}}b_{\text{},\text{},\text{}}^{},\gamma _{\mathrm{𝟐},\mathrm{𝟐}}^{}=_{\text{df}}b_{\text{},\text{},\text{}}^{},$$ and by using ($`\gamma `$1), (hex 1), (hex 1a), (hex 2) and (hex 2a) as clauses in an inductive definition. The equations of $`𝚪`$ certainly hold in Â for this defined structure because Â is a preorder, as we said above. To finish the proof that $`𝚪`$ and Â are isomorphic categories, it remains only to check that the clauses of the inductive definitions of $`_n`$, $`\gamma _{A,B}^{}`$ and $`\gamma _{A,B}^{}`$ hold as equations in $`𝚪`$ for $`b_{A,B,C}^{}`$, $`b_{A,B,C}^{}`$ and $``$ defined as they are defined in $`𝚪`$. This is done by using essentially (assoc 1$``$) and (assoc 2$``$). So $`𝚪`$ is isomorphic to Â, and is hence a preorder. If we have instead of Â the free monoidal category without unit, i.e. the free associative category, $`\text{}^{}`$ generated by an arbitrary nonempty set of objects $`𝒫`$, conceived as a discrete category, then, instead of $`𝚪`$, the analogous category $`𝚪^{}`$ isomorphic to $`\text{}^{}`$ would have as generators $`𝒫\{\mathrm{𝟐}\}`$. Every object of $`𝚪^{}`$ different from a member of $`𝒫`$ can be written in the form $$(\mathrm{}(C_np_n)\mathrm{}_2p_2)_1p_1$$ for $`C`$ an object of $`𝚪`$ (more precisely, a member of $`_2`$), $`n=\|C\|`$ and $`p_1,\mathrm{},p_n`$ members of $`𝒫`$. For the arrows $`\gamma _{A,B}^{}`$ and $`\gamma _{A,B}^{}`$ we would assume that $`p_{\|A\|}`$ in $`A`$ coincides with $`p_1`$ in $`B`$, and the equations ($`\gamma `$1) and (unit $``$) would have to be adapted. We have seen above how Mac Lane’s pentagon arises from a Yang-Baxter hexagon by collapsing, according to (assoc 2), the vertices corresponding to $`(\mathrm{𝟐}_1\mathrm{𝟐})_3\mathrm{𝟐}`$, i.e. $`BAC`$, and $`(\mathrm{𝟐}_2\mathrm{𝟐})_1\mathrm{𝟐}`$, i.e. $`BCA`$, into a single vertex corresponding to $`(\text{}\text{})(\text{}\text{})`$. We can apply this collapsing procedure based on (assoc 2) to the three-dimensional permutohedron (whose vertices correspond to permutations of four letters and edges to transpositions of adjacent letters) in order to obtain the three-dimensional associahedron (whose vertices correspond to planar binary trees with five leaves and edges to arrow terms of Â with a single $`b^{}`$), and afterwards we can proceed to higher dimensions. The function that corresponds to our procedure is described in . Our paper provides a motivation for that function.
warning/0506/hep-ex0506039.html	ar5iv	text	# New resonances in B-meson decays ## 1 Introduction $`B`$-mesons have proved to be a rich source of new particles. The B-factories at KEK (Belle) and SLAC (BaBar) extensively use the exclusive production of $`J/\psi `$, $`\chi _{c1}`$ and $`\eta _c`$ charmonia for the CP violation measurements. In addition to these conventional charmonia Belle observed the production of $`\eta _c(2S)`$ $`^\mathrm{?}`$ and $`\psi (3770)`$ $`^\mathrm{?}`$ in exclusive $`B(c\overline{c})K^+`$ decays. Belle and BaBar also observed inclusive $`\chi _{c2}`$ production in $`B`$ decays $`^\mathrm{?}`$. In 2003, by analyzing the $`B^+J/\psi \pi ^+\pi ^{}K^+`$ decays, Belle observed a narrow charmonium-like new state (denoted as $`X(3872)`$) decaying into $`J/\psi \pi ^+\pi ^{}`$ $`^\mathrm{?}`$. Recently Belle reported the observation of $`Y(3940)J/\psi \omega `$ in $`B^+Y(3940)K^+`$ decays $`^\mathrm{?}`$. All observed resonances – from $`\eta _c(2S)`$ to $`Y(3940)`$ – are difficult to reconstruct without the constraints provided by $`B`$ decays, as they decay to high-multiplicity final states. $`B`$ mesons also provide an excellent opportunity to test different hypotheses for the $`J^P`$ quantum numbers of these resonances via decay angle analysis. ## 2 Observation of $`X(3872)`$ and it’s properties Just after the discovery of $`X(3872)`$ by Belle, this new state was confirmed by the CDF, D0 and BaBar collaborations $`^\mathrm{?}`$. Its mass was measured to be $`3871.9\pm 0.5`$ MeV which is very close to the $`D^0\overline{D}^0`$ threshold of $`3871.3\pm 1.0`$ MeV. All these first measurements directly provide a lot of information on the $`X(3872)`$ properties $`^\mathrm{?}`$: * the two pion mass from the $`X(3872)J/\psi \pi ^+\pi ^{}`$ decay tends to peak near the $`\rho ^0`$ mass, consistent with positive C-parity of the $`X(3872)`$; * the decays $`X(3872)\chi _{c1,2}\gamma `$ are not seen; this likely excludes the $`1^3D_{2,3}`$ ($`\psi _{2,3}`$) assignment for the $`X(3872)`$; * the decay $`X(3872)D\overline{D}`$ is suppressed or forbidden $`^\mathrm{?}`$; this, together with its narrow width ($`\mathrm{\Gamma }<2.7`$ MeV), suggests $`J^P=0^+,1^{},2^+,\mathrm{}`$ are ruled out; CDF and D0 have measured the $`X(3872)`$ production properties to be very similar to those of the $`\psi (2S)`$ $`^\mathrm{?}`$. BaBar reported a null search for the charged $`X(3872)`$ partners in $`BJ/\psi \pi ^+\pi ^0K^+`$ decays $`^\mathrm{?}`$ and a null search for $`X(3872)J/\psi \eta `$ decay $`^\mathrm{?}`$. The former rules out the isovector hypothesis and the latter excludes the presence of gluonic degrees of freedom in the $`X(3872)`$ wave function. Fig. 1 shows the $`M(\pi ^+\pi ^{})`$ spectrum from the updated analysis of $`X(3872)J/\psi \pi ^+\pi ^{}`$ decays by Belle (253 fb<sup>-1</sup>). The $`\rho ^0`$ signal is strong and supports $`C(X(3872))=+1`$. Recently Belle has found evidence for another decay mode of $`X(3872)`$: $`X(3872)J/\psi \omega ^{}`$ where $`\omega ^{}`$ is virtual and reconstructed in the $`\pi ^+\pi ^{}\pi ^0`$ channel. According to Swanson’s model $`^\mathrm{?}`$, this observation supports the $`D^0\overline{D}^0`$ molecular interpretation of $`X(3872)`$. The full set of the most recent Belle results on the $`X(3872)`$ properties can be found elsewhere $`^\mathrm{?}`$. ## 3 Observation of $`Y(3940)`$ By analyzing exclusive $`B^+J/\psi \pi ^+\pi ^{}\pi ^0K^+`$ decays Belle observed a new resonance $`Y(3940)`$ decaying to $`J/\psi \omega `$ $`^\mathrm{?}`$. Fig. 2 shows the $`M(J/\psi \omega )`$ distribution for the $`B`$-meson candidates. The curve in Fig. 2 (a) indicates the result of a fit with only a phase-space-like two-body threshold function. The curve in Fig. 2 (b) shows the result of a fit that includes an S-wave Breit-Wigner resonance term. The mass and width of $`Y(3940)`$ were measured to be $`3943\pm 11\pm 13`$ MeV and $`87\pm 22\pm 26`$ MeV, respectively. The observed state is above the $`D\overline{D}^{()}`$ threshold and would decay predominantly to $`D\overline{D}`$ and/or $`D\overline{D}^{}`$ if it is a $`c\overline{c}`$ charmonium. In contrast, for a $`c\overline{c}gluon`$ hybrid the open charm decay modes are suppressed or forbidden. So the observed $`Y(3940)`$ is a possible candidate for the first $`c\overline{c}gluon`$ hybrid state. ## 4 Search for $`B^+h_cK^+`$ The $`1^1P_1`$, $`J^P=1^+`$ $`h_c`$ has for a long time been a missing state. Recently CLEO reported the observation of $`h_c`$ in $`\psi (2S)h_c\pi ^0`$, $`h_c\eta _c\gamma `$ decays $`^\mathrm{?}`$. The mass was measured to be $`3524.4\pm 0.6\pm 0.4`$ MeV - in agreement with theoretical expectations that the $`M(h_c)`$ is close to the c.o.g. of of the $`<1^3P_J>`$ triplet. Belle searched for $`h_c`$ production in exclusive $`B^+h_cK^+`$, $`h_c\eta _c\gamma `$ decays. Fig. 3 shows the $`M(\eta _c\gamma )`$ for the $`B`$-candidates. No evidence for a signal around the CLEO $`h_c`$ mass is seen. Belle set an upper limit for $`(B^+h_cK^+)\times (h_c\eta _c\gamma )`$ that is less than $`1.5\times 10^4`$ for $`M(h_c)3520`$ MeV $`^\mathrm{?}`$. ## 5 Study of $`D_{sJ}(2317)`$ and $`D_{sJ}(2460)`$ The $`D_{sJ}(2317)^+`$ and $`D_{sJ}(2460)^+`$ mesons were observed by BaBar $`^\mathrm{?}`$ and CLEO $`^\mathrm{?}`$. Belle confirmed these states and observed their production in exclusive $`BD_{sJ}^+\overline{D}^{()}`$ decays $`^\mathrm{?}`$. The observation of these decays allowed us to perform decay angle analysis. Belle data support $`J=0`$ for $`D_{sJ}(2317)`$ and $`J=1`$ for $`D_{sJ}(2460)`$ $`^\mathrm{?}`$. ## 6 Summary $`B`$ mesons provide a clean environment for the observation of yet-unseen charmonia and other new unexpected resonances and the understanding of their properties. The nature of new resonances $`X(3872)`$ and $`Y(3940)`$ remains unclear so far. They could be either $`c\overline{c}`$ states or exotic hadrons: $`D^0\overline{D}^0`$ molecular ($`X(3872)`$) and $`c\overline{c}gluon`$ hybrid ($`Y(3940)`$). ## Acknowledgments We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the NII for valuable computing and Super-SINET network support. We acknowledge support from MEXT and JSPS (Japan); ARC and DEST (Australia); NSFC (contract No. 10175071, China); DST (India); the BK21 program of MOEHRD and the CHEP SRC program of KOSEF (Korea); KBN (contract No. 2P03B 01324, Poland); MIST (Russia); MHEST (Slovenia); SNSF (Switzerland); NSC and MOE (Taiwan); and DOE (USA). ## References
warning/0506/math0506345.html	ar5iv	text	# Elementary proofs of Paley–Wiener theorems for the Dunkl transform on the real line ## 1. Introduction and overview The Paley–Wiener theorem for the Dunkl transform $`D_k`$ with multiplicity $`k`$ (where $`Rek0`$) on the real line states that a smooth function $`f`$ has support in the bounded interval $`[R,R]`$ if, and only if, its transform $`D_kf`$ is an entire function which satisfies the usual growth estimates as they are required in the (special) case of the Fourier transform. Various proofs of this result are known, all of which use explicit formulas available in this one-dimensional setting (see Remark 6 for more details). In this paper, however, we present an alternative proof which does not use such explicit expressions, being based almost solely on the formal properties of the transform. Along the same lines, we also obtain a Paley–Wiener theorem for $`L^2`$-functions for $`k0`$. The case $`k=0`$ specializes to the Fourier transform, and to our knowledge the proofs of both Paley–Wiener theorems are new even in this case. In addition, we establish two identities in the spirit of Bang \[4, Theorem 1\]. These results could be called real Paley–Wiener theorems (although terminology is not yet well-established), since they relate certain growth rates of a function on the real line to the support of its transform. The approach at this point is inspired by similar techniques in . Our results in this direction partially overlap with , but the new proofs are considerably simpler, as they are again based almost solely on the formal properties of the transform. We will comment on this in more detail later on, as these results will have been established. For $`k=0`$, one retrieves Bang’s result; we feel that the present method of proof, which, e.g., does not use the Paley–Wiener theorem for smooth functions, but rather implies it, is then more direct than that in . The rather unspecific and formal structure of the proofs suggests that the methods can perhaps be put to good use for other integral transforms with a symmetric kernel, both for the Paley–Wiener theorems and the equalities in the spirit of Bang (cf. Remark 6). The structure of the proof is also such that, if certain combinatorial problems can be surmounted, a proof of the Paley–Wiener theorem for the Dunkl transform for invariant balanced compact convex sets in arbitrary dimension might be possible. This would be further evidence for the validity of this theorem for invariant compact convex sets (cf. \[13, Conjecture 4.1\]), but at the time of writing this higher-dimensional result has not been established. This paper is organized as follows. In Section 2 the necessary notations and previous results are given. Section 3 contains the Paley–Wiener theorem for smooth functions and can—perhaps—serve as a model for a proof of such a theorem in other contexts. The rest of the paper is independent of this section. Section 4 is concerned with the real Paley–Wiener theorem in the $`L^p`$-case and the $`L^2`$-case is settled in Section 5. ## 2. Dunkl operators and the Dunkl transform on $``$ The Dunkl operators and the Dunkl transform were introduced for arbitrary root systems by Dunkl . In this section we recall some basic properties for the one-dimensional case of $`A_1`$, referring to \[17, Sections 1 and 2\] for a more comprehensive overview and to for details. We suppress the various explicit formulas which are known in this one-dimensional context (as these are not necessary for the proofs), thus emphasizing the basic structure of the problem which might lead to generalizations to the case of arbitrary root systems. Let $`k`$, and consider the Dunkl operator $`T_k`$ $$T_kf(x)=f^{}(x)+k\frac{f(x)f(x)}{x}(fC^{\mathrm{}}(),x).$$ Rewriting this as (1) $$T_kf(x)=f^{}(x)+k_1^1f^{}(tx)𝑑t,$$ it follows that $`T_k`$ maps $`C^{\mathrm{}}()`$, $`C_c^{\mathrm{}}()`$ and the Schwartz space $`𝒮()`$ into themselves. If $`Rek0`$, as we will assume for the remainder of this section, then, for each $`\lambda `$, there exists a unique holomorphic solution $`\psi _\lambda ^k:`$ of the differential-reflection problem (2) $$\{\begin{array}{cc}T_kf=i\lambda f,\hfill & \\ f(0)=1.\hfill & \end{array}$$ The map $`(z,\lambda )\psi _\lambda ^k(z)`$ is entire on $`^2`$, and we have the estimate (3) $$\left\|\psi _\lambda ^k(z)\right\|e^{\|Im\lambda z\|}(\lambda ,z).$$ In view of (3) the Dunkl transform $`D_kf`$ of $`fL^1(,\|w_k(x)\|dx)`$, where the complex-valued weight function $`w_k`$ is given by $`w_k(x)=\|x\|^{2k}`$, is meaningfully defined by (4) $$D_kf(\lambda )=\frac{1}{c_k}_{}f(x)\psi _\lambda ^k(x)w_k(x)𝑑x(\lambda ),$$ where $$c_k=_{}e^{\frac{\|x\|^2}{2}}w_k(x)𝑑x0.$$ We note that $`D_0`$ is the Fourier transform on $``$. From (3) we conclude that $`D_kf`$ is bounded for such $`f`$, in fact (5) $$\|D_kf(\lambda )\|\frac{1}{\|c_k\|}_{}\|f(x)\|\|w_k(x)\|𝑑x(\lambda ,fL^1(,\|w_k(x)\|dx)).$$ The Dunkl transform is a topological isomorphism of $`𝒮()`$ onto itself, the inverse transform $`D_k^1`$ being given by $$D_k^1f(x)=\frac{1}{c_k}_{}f(\lambda )\psi _\lambda ^k(x)w_k(\lambda )𝑑\lambda =D_kf(x)(f𝒮(),x).$$ The operator $`T_k`$ is anti-symmetric with respect to the weight function $`w_k`$, i.e., (6) $$T_kf,g_k=f,T_kg_k,$$ for $`f𝒮()`$ and $`gC^{\mathrm{}}()`$ such that both $`g`$ and $`T_kg`$ are of at most polynomial growth. Here $`f,g_k`$ is defined by $$f,g_k=_{}f(x)g(x)w_k(x)𝑑x,$$ for functions $`f`$ and $`g`$ such that $`fgL^1(,\|w_k(x)\|dx)`$. In particular, (6) yields the intertwining identity $$D_k(T_kf)(\lambda )=i\lambda (D_kf)(\lambda )(f𝒮(),\lambda ).$$ Furthermore, for $`\lambda ,z,s`$ the symmetry properties $`\psi _\lambda ^k(z)=\psi _z^k(\lambda )`$ and $`\psi _{s\lambda }^k(z)=\psi _\lambda ^k(sz)`$ are valid. Using the first of these and (5), an application of Fubini gives (7) $$D_kf,g_k=f,D_kg_k(f,gL^1(,\|w_k(x)\|dx)).$$ If $`k0`$, the Plancherel theorem states that $`D_k`$ preserves the weighted two-norm on $`L^1(,w_k(x)dx)L^2(,w_k(x)dx)`$ and extends to a unitary operator on $`L^2(,w_k(x)dx)`$. ## 3. Paley–Wiener theorem for smooth functions The method of proof in this section is inspired by results of Bang . To be more specific, for $`R>0`$ let $`_R()`$ denote the space of entire functions $`f`$ with the property that, for all $`n\{0\}`$, there exists a constant $`C_{n,f}>0`$ such that $$\|f(z)\|C_{n,f}(1+\|z\|)^ne^{R\|Imz\|}(z).$$ Then, if $`k0`$ and $`f_R()`$, we will establish that (cf. ) (8) $$sup\{\|\lambda \|:\lambda \mathrm{supp}D_kf\}\underset{n\mathrm{}}{lim\; inf}T_k^nf_{\mathrm{}}^{1/n}\underset{n\mathrm{}}{lim\; sup}T_k^nf_{\mathrm{}}^{1/n}R,$$ after which the proof of the Paley–Wiener theorem for smooth functions is a mere formality. Starting towards the third of these inequalities, we first use (1) to gain control over repeated Dunkl derivatives. ###### Lemma 1. Let $`k`$, $`fC^{\mathrm{}}()`$ and $`n`$. Then $$T_k^nf(x)=(T_k^{n1}(f^{}))(x)+k_1^1t^{n1}(T_k^{n1}(f^{}))(tx)𝑑t(x).$$ ###### Proof. For $`gC^{\mathrm{}}()`$ and $`m\{0\}`$, let $`I_{g,m}C^{\mathrm{}}()`$ be defined by $$I_{g,m}(x)=_1^1t^mg(tx)𝑑t(x).$$ Using (1), we then find $`(T_kI_{g,m})(x)`$ $`={\displaystyle _1^1}t_1^{m+1}g^{}(t_1x)𝑑t_1+k{\displaystyle _1^1}{\displaystyle _1^1}t_1^{m+1}g^{}(t_1t_2x)𝑑t_1𝑑t_2`$ $`={\displaystyle _1^1}t_1^{m+1}\left(g^{}(t_1x)+k{\displaystyle _1^1}g^{}(t_1t_2x)𝑑t_2\right)𝑑t_1`$ $`={\displaystyle _1^1}t_1^{m+1}(T_kg)(t_1x)𝑑t_1=I_{T_kg,m+1}(x).`$ We conclude that $`T_kI_{g,m}=I_{T_kg,m+1}`$. Since (1) can be written as $`T_kf=f^{}+kI_{f^{},0}`$ one has $$T_k^nf=T_k^{n1}(f^{}+kI_{f^{},0})=T_k^{n1}(f^{})+kI_{T_k^{n1}(f^{}),n1},$$ which is the statement in the lemma. ∎ It follows from Lemma 1 that $$\|T_k^nf(x)\|\left(1+\frac{2\|k\|}{n}\right)\underset{y[\|x\|,\|x\|]}{sup}\|(T_k^{n1}(f^{}))(y)\|(x),$$ and induction then yields the following basic estimate, which is more explicit than \[7, Prop. 2.1\]. ###### Corollary 2. Let $`k`$, $`fC^{\mathrm{}}()`$ and $`n`$. Then $$\|T_k^nf(x)\|\frac{\mathrm{\Gamma }(n+1+2\|k\|)}{n!\mathrm{\Gamma }(1+2\|k\|)}\underset{y[\|x\|,\|x\|]}{sup}\|f^{(n)}(y)\|(x).$$ The third inequality in (8) can now be settled. ###### Proposition 3. Let $`k,R>0`$, and suppose $`f:`$ is an entire function such that $$\|f(z)\|Ce^{R\|Imz\|}(z),$$ for some positive constant $`C`$. Then, for all $`n`$, $`T_k^nf`$ is bounded on the real line, and $$\underset{n\mathrm{}}{lim\; sup}T_k^nf_{\mathrm{}}^{1/n}R.$$ ###### Proof. We have, for any $`r>0`$, $$f^{(n)}(z)=\frac{n!}{2\pi i}_{\|\zeta z\|=r}\frac{f(\zeta )}{(\zeta z)^{n+1}}𝑑\zeta (z).$$ If $`\|\zeta z\|=r`$, then $$\|f(\zeta )\|Ce^{R(\|Imz\|+r)},$$ implying $$\|f^{(n)}(z)\|C\frac{n!}{r^n}e^{R(\|Imz\|+r)}(z).$$ Choosing $`r=n/R`$, so that $`r>0`$ if $`n`$, we find $$\|f^{(n)}(z)\|C\frac{n!e^n}{n^n}R^ne^{R\|Imz\|}(n,z),$$ whence $`f^{(n)}_{\mathrm{}}Cn!e^nn^nR^n`$, for $`n`$. Combining this with Corollary 2 yields $$T_k^nf_{\mathrm{}}C\frac{e^n\mathrm{\Gamma }(n+1+2\|k\|)}{n^n\mathrm{\Gamma }(1+2\|k\|)}R^n(n).$$ The result now follows from Stirling’s formula. ∎ As to the first inequality in (8), it is actually easy to prove that it holds for the norm $`_{k,p}`$ in $`L^p(,w_k(x)dx)`$ for arbitrary $`1p\mathrm{}`$ (not just for $`p=\mathrm{}`$), as is shown by the following lemma. It should be noted that this result can be generalized—with different proofs—to complex multiplicities (cf. Lemma 7) and to $`L^p`$-functions for $`k0`$ (cf. Theorem 10), but we present it here separately nevertheless, in order to illustrate that for the case $`k0`$, the proof of one of the crucial inequalities (as far as the Paley–Wiener theorem is concerned) is rather elementary and intuitive. ###### Lemma 4. Let $`k0`$, $`1p\mathrm{}`$ and $`f𝒮()`$. Then in the extended positive real numbers, (9) $$\underset{n\mathrm{}}{lim\; inf}T_k^nf_{k,p}^{1/n}sup\{\|\lambda \|:\lambda \mathrm{supp}D_kf\}.$$ ###### Proof. Suppose $`0\lambda _0\mathrm{supp}D_kf`$ and let $`0<\epsilon <\|\lambda _0\|`$. Define $$\varphi (\lambda )=\overline{D_kf(\lambda )}(\lambda ).$$ Then, with $`q`$ denoting the conjugate exponent and using (7), we find $`T_k^{2n}f_{k,p}D_k\varphi _{k,q}`$ $`\|T_k^{2n}f,D_k\varphi _k\|=\|D_k(T_k^{2n}f),\varphi _k\|`$ $`=\left\|{\displaystyle _{}}(i\lambda )^{2n}D_kf(\lambda )\varphi (\lambda )w_k(\lambda )𝑑\lambda \right\|`$ $`={\displaystyle _{}}\lambda ^{2n}\|D_kf(\lambda )\|^2w_k(\lambda )𝑑\lambda `$ $`(\|\lambda _0\|\epsilon )^{2n}{\displaystyle _{\|\lambda \|\|\lambda _0\|\epsilon }}\|D_kf(\lambda )\|^2w_k(\lambda )𝑑\lambda .`$ With $`\psi (\lambda )=\overline{D_k(T_kf)(\lambda )}`$ we similarly get $$T_k^{2n+1}f_{k,p}D_k\psi _{k,q}(\|\lambda _0\|\epsilon )^{2n+1}_{\|\lambda \|\|\lambda _0\|\epsilon }\|\lambda \|\|D_kf(\lambda )\|^2w_k(\lambda )𝑑\lambda .$$ These two estimates together yield $$\underset{n\mathrm{}}{lim\; inf}T_k^nf_{k,p}^{1/n}\|\lambda _0\|\epsilon ,$$ and the lemma follows. ∎ Now that (8) has been established, we come to the Paley–Wiener theorem for smooth functions. Introducing notation, for $`R>0`$, we let $`C_R^{\mathrm{}}()`$ denote the space of smooth functions on $``$ with support in $`[R,R]`$. Its counterpart under the Dunkl transform, the space $`_R()`$, was defined at the beginning of this section. ###### Theorem 5 (Paley–Wiener theorem for smooth functions). Let $`R>0`$ and $`k0`$. Then the Dunkl transform $`D_k`$ is a bijection from $`C_R^{\mathrm{}}()`$ onto $`_R()`$. ###### Proof. If $`fC_R^{\mathrm{}}()`$, then it is easy to see that $`D_kf_R()`$ \[12, Corollary 4.10\]. Now assume that $`f_R()`$. Using Cauchy’s integral representation as in the proof of Proposition 3, we retrieve the well-known fact that $`f𝒮()`$. From (8) we infer that $`D_kf`$ has support in $`[R,R]`$. Since $`D_k^1f(x)=D_kf(x)`$, for $`x`$, the same is true for $`D_k^1f`$, as was to be proved. ∎ ###### Remark 6. 1. By holomorphic continuation and continuity, cf. , one sees that Theorem 5 also holds in the more general case $`Rek0`$. Alternatively, one can use Lemma 7 below instead of Lemma 4, which establishes (8) also in the case $`Rek0`$, and then the above direct proof is again valid. 2. We emphasize that the present proof does not use any explicit formulas for the Dunkl kernel in one dimension, contrary to the alternative methods of proof in (where Weyl fractional integral operators are used), (where asymptotic results for Bessel functions are needed) and (where various integral operators, Dunkl’s intertwining operator and the Paley–Wiener theorem for the Fourier transform all play a role). Also, the fact that a contour shifting argument for the transform is usually not possible (since $`w_k`$ generically has no entire extension) is no obstruction. Given this unspecific nature, it is possible that the present method can be applied to other transforms as well, although the symmetry of the kernel—as reflected in (7), which was used in the proof of Lemma 4 and which will again be used in the proof of the alternative Lemma 7 below—is perhaps necessary. The same suggestion applies to the results in the remaining sections of this paper. ## 4. Real Paley–Wiener theorem for $`L^p`$-functions We will now consider the real Paley–Wiener theorem for $`L^p`$-functions in the spirit of Bang . The result is first proved for Schwartz functions in Theorem 8 and subsequently for the general case in Theorem 10. Let $`_{Rek,p}`$ denote the $`L^p(,\|w_k(x)\|dx)`$-norm, for $`1p\mathrm{}`$. Then we have the following generalization of Lemma 4 to complex multiplicities. ###### Lemma 7. Let $`Rek0`$, $`1p\mathrm{}`$ and $`f𝒮()`$. Then in the extended positive real numbers, (10) $$\underset{n\mathrm{}}{lim\; inf}T_k^nf_{Rek,p}^{1/n}sup\{\|\lambda \|:\lambda \mathrm{supp}D_kf\}.$$ ###### Proof. Let $`0\lambda _0\mathrm{supp}D_kf`$ and choose $`ϵ>0`$ such that $`0<2\epsilon <\|\lambda _0\|`$. Also choose $`\varphi C_c^{\mathrm{}}()`$ such that $`\mathrm{supp}\varphi [\lambda _0\epsilon ,\lambda _0+\epsilon ]`$, and $`D_kf,\varphi _k0`$. Define $`\varphi _n(\lambda )=\lambda ^n\varphi (\lambda )`$ and $`P_n(x)=x^n`$ for $`n\{0\}`$. Then $$(1+P_N(x))(D_k\varphi _n)(x)=\frac{1}{c_k}_{\lambda _0\epsilon }^{\lambda _0+\epsilon }\left(1+(iT_k)^N\right)(\lambda ^n\varphi (\lambda ))\psi _x^k(\lambda )w_k(\lambda )𝑑\lambda (N\{0\}).$$ We fix $`N`$ such that $`N`$ is even and $`N>2Rek+1`$. Corollary 2 and the binomial formula imply that $$\left\|\left(1+(iT_k)^N\right)(\lambda ^n\varphi (\lambda ))\right\|C_1n^N(\|\lambda _0\|\epsilon )^n(n\{0\},\lambda ),$$ where $`C_1`$ is a positive constant. This yields the estimates $`D_k\varphi _n_{Rek,q}`$ $`(1+P_N)^1_{Rek,q}(1+P_N)D_k\varphi _n_{\mathrm{}}`$ $`{\displaystyle \frac{2\epsilon }{\|c_k\|}}C_1n^N\left(\|\lambda _0\|\epsilon \right)^n(1+P_N)^1_{Rek,q}`$ $`C_2n^N\left(\|\lambda _0\|\epsilon \right)^n`$ for all $`n>N`$, where $`C_2`$ is a positive constant and $`q`$ is the conjugate exponent. Using (7), the identity $`D_k(P_n\varphi _n)=(i)^nT_k^nD_k\varphi _n`$ and Hölder’s inequality, we therefore get $`\|D_kf,\varphi _k\|`$ $`=\|D_kf,P_n\varphi _n_k\|=\|f,D_k(P_n\varphi _n)_k\|=\|f,T_k^n(D_k\varphi _n)_k\|`$ $`=\left\|T_k^nf,D_k\varphi _n_k\right\|T_k^nf_{Rek,p}D_k\varphi _n_{Rek,q}`$ $`C_2n^N\left(\|\lambda _0\|\epsilon \right)^nT_k^nf_{Rek,p},`$ whence $$\underset{n\mathrm{}}{lim\; inf}T_k^nf_{Rek,p}^{1/n}\underset{n\mathrm{}}{lim\; inf}(C_2n^N)^{1/n}\left(\|\lambda _0\|\epsilon \right)\|D_kf,\varphi _k\|^{1/n}=\left(\|\lambda _0\|\epsilon \right),$$ establishing the lemma. ∎ Using the Paley–Wiener theorem, we can extend the inequality in Proposition 3 to the norms $`_{Rek,p}`$, $`1p\mathrm{}`$, for $`Rek0`$, and we thus have the following theorem. ###### Theorem 8 (real Paley–Wiener theorem for Schwartz functions). Let $`Rek0`$, $`1p\mathrm{}`$ and $`f𝒮()`$. Then in the extended positive real numbers, $$\underset{n\mathrm{}}{lim}T_k^nf_{Rek,p}^{1/n}=sup\{\|\lambda \|:\lambda \mathrm{supp}D_kf\}.$$ ###### Proof. In view of Lemma 7 it only remains to be shown that $$\underset{n\mathrm{}}{lim\; sup}T_k^nf_{Rek,p}^{1/n}R,$$ if $`f𝒮()`$ is such that $`\mathrm{supp}D_kf[R,R]`$ for some finite $`R>0`$. Using the inversion formula and the intertwining properties of the transform we have (11) $$x^NT_k^nf(x)=\frac{i^{N+n}}{c_k}_R^RT_k^N(P_nD_kf)(\lambda )\psi _\lambda ^k(x)w_k(\lambda )𝑑\lambda (n,N\{0\}),$$ where again $`P_n(x)=x^n`$. Now Corollary 2 and the binomial formula imply that $$T_k^N(P_nD_kf)_{\mathrm{}}C_1n^NR^n,$$ where $`C_1`$ is a constant depending on $`f`$ and $`N`$. Therefore (11) yields that (12) $$(1+P_N)T_k^nf_{\mathrm{}}C_2n^NR^{n+1},$$ where $`C_2`$ is again a constant depending on $`f`$ and $`N`$. We fix $`N`$ such that $`N`$ is even and $`N>2Rek+1`$. Then the observation $$T_k^nf_{Rek,p}(1+P_N)^1_{Rek,p}(1+P_N)T_k^nf_{\mathrm{}}$$ and (12) establish the result. ∎ ###### Remark 9. Theorem 8 is new for complex $`k`$. For real $`k`$, the result can be found in , where it is proved using the Plancherel theorem for the Dunkl transform, the Riesz–Thorin convexity theorem and the theory of Sobolev spaces for Dunkl operators. For $`k0`$, we will now generalize Theorem 8 to the $`L^p`$-case in Theorem 10, using the structure of $`𝒮()`$ as an associative algebra under the Dunkl convolution $`_k`$. We refer to for details on this subject. Let $`k0`$, and define a distributional Dunkl transform $`D_k^df`$ of $`fL^p(,w_k(x)dx)`$ by transposition $$D_k^df,\varphi =f,D_k\varphi _k(\varphi 𝒮()).$$ Clearly $`D_k^d`$ is injective on $`L^p(,w_k(x)dx)`$. Furthermore, from $$\|D_k^df,\varphi \|f_pD_k\varphi _{\mathrm{}}(\varphi 𝒮()),$$ we see that $`D_k^df`$ is a tempered distribution. Note also that for $`fL^p(,w_k(x)dx)`$ and $`1p2`$ we have $$D_k^df,\varphi =D_kf,\varphi _k(\varphi 𝒮()).$$ by (7) and density of $`𝒮()`$ in $`L^p(,w_k(x)dx)`$. Thus, for $`fL^p(,w_k(x)dx)`$ and $`1p2`$, $`D_k^df`$ as defined above corresponds to the distribution $`(D_kf)w_k`$, implying that in this case $`\mathrm{supp}D_k^df=\mathrm{supp}D_kf`$. ###### Theorem 10 (real Paley–Wiener theorem for $`L^p`$-functions). Let $`k0`$, $`1p\mathrm{}`$ and $`fC^{\mathrm{}}()`$ be such that $`T_k^nfL^p(,w_k(x)dx)`$, for all $`n\{0\}`$. Then in the extended positive real numbers, $$\underset{n\mathrm{}}{lim}T_k^nf_{k,p}^{1/n}=sup\{\|\lambda \|:\lambda \mathrm{supp}D_k^df\}.$$ If in addition $`fL^s(,w_k(x)dx)`$, for some $`1s2`$, then the distribution $`D_k^df`$ corresponds to the function $`(D_kf)w_k`$, and the support of $`D_k^df`$ in the right hand side is equal to the support of $`D_kf`$ as a distribution. ###### Proof. First we note that (10) also holds for $`f`$ as above: in the proof of Lemma 7 we just have to change $`D_kf,_k`$ into $`D_k^df,`$. Therefore, it only remains to be shown that (13) $$\underset{n\mathrm{}}{lim\; sup}T_k^nf_{k,p}^{1/n}R,$$ if $`f`$ is as in the theorem and such that $`\mathrm{supp}D_k^df[R,R]`$ for some finite $`R>0`$. To this end, choose $`\epsilon >0`$, and fix a function $`\varphi _\epsilon 𝒮()`$ such that $`D_k^1\varphi _\epsilon =1`$ on $`[R,R]`$ and $`D_k^1\varphi _\epsilon =0`$ outside $`[R\epsilon ,R+\epsilon ]`$. We have from \[19, Proposition 3\] that $`D_k^1(\varphi _k\psi )=(D_k^1\varphi )(D_k^1\psi )`$, for all $`\varphi ,\psi 𝒮()`$. With $`P_n(x)=x^n`$, we thus find for arbitrary $`\varphi 𝒮()`$ that $`f,\varphi _kT_k^n\varphi _\epsilon _k`$ $`=D_k^df,D_k^1(\varphi _kT_k^n\varphi _\epsilon )=(i)^nD_k^df,P_n(D_k^1\varphi )(D_k^1\varphi _\epsilon )`$ $`=(i)^nD_k^df,P_nD_k^1\varphi =f,T_k^n\varphi _k.`$ Furthermore, from loc.cit., one knows that $`\varphi _k\psi _{k,q}4\varphi _{k,q}\psi _{k,1}`$, for all $`\varphi ,\psi 𝒮()`$, where $`q`$ is the conjugate exponent of $`p`$. If we combine these two results with (6) and Hölder’s reverse inequality, we infer that $$T_k^nf_{k,p}=\underset{\varphi }{sup}\|T_k^nf,\varphi _k\|=\underset{\varphi }{sup}\|f,T_k^n\varphi _k\|=\underset{\varphi }{sup}\|f,\varphi _kT_k^n\varphi _\epsilon _k\|4f_{k,p}T_k^n\varphi _\epsilon _{k,1},$$ where the supremum is over all functions $`\varphi 𝒮()`$ with $`\varphi _{k,q}=1`$. From Theorem 8 we therefore conclude that $$\underset{n\mathrm{}}{lim\; sup}T_k^nf_{k,p}^{1/n}R+\epsilon ,$$ proving (13). The statement on supports was established in the discussion preceding the theorem. ∎ ###### Remark 11. For even functions, when the Dunkl transform reduces to the Hankel transform, the previous result can already be found in . Also using Dunkl convolution, and closely following the approach in , Theorem 10 has previously been established in . Our proof is considerably shorter than the proof in loc.cit. ## 5. Paley–Wiener theorem for $`L^2`$-functions (for $`k0`$) In this section we assume that $`k0`$. For $`R>0`$ and $`1p\mathrm{}`$, we define $`L_R^p(,w_k(x)dx)`$ to be the subspace of $`L^p(,w_k(x)dx)`$ consisting of those functions with distributional support in $`[R,R]`$, and we let $`_R^{p,k}()`$ denote the space of entire functions $`f:`$ which belong to $`L^p(,w_k(x)dx)`$ when restricted to the real line and which are such that $$\|f(z)\|C_fe^{R\|Imz\|}(z),$$ for some positive constant $`C_f`$. ###### Theorem 12 (Paley–Wiener theorem for $`L^2`$-functions). Let $`R>0`$ and $`k0`$. Then the Dunkl transform $`D_k`$ is a bijection from $`L_R^2(,w_k(x)dx)`$ onto $`_R^{2,k}()`$. ###### Proof. Let $`fL_R^2(,w_k(x)dx)`$. Then $`fL^1(,w_k(x)dx)`$, and (3) and (4) together with the Plancherel theorem imply that $`D_kf_R^{2,k}()`$. Conversely, let $`f_R^{2,k}()`$. By the Plancherel theorem one has $`D_k^1fL^2(,w_k(x)dx)`$. In addition, Proposition 3 and Theorem 10 (with $`p=\mathrm{}`$) show that $`\mathrm{supp}D_kf[R,R]`$. The same is then true for $`D_k^1f`$, and the result follows. ∎ ###### Remark 13. 1. If $`f_R^{p,k}()`$ and $`1p\mathrm{}`$, then using Proposition 3 and Theorem 10 as in the proof of Theorem 12, one sees that $`D_k^df`$ has support in $`[R,R]`$. In particular, for $`1p2`$ with conjugate exponent $`q`$, we conclude that $`D_k`$ maps $`_R^{p,k}()`$ into $`L_R^q(,w_k(x)dx)`$. For even functions, when the Dunkl transform can be identified with the Hankel transform, the latter result can be found in . 2. As an ingredient for the discussion of the relation with the literature on the Fourier transform, let us make the preliminary observation that, for $`R>0`$ and $`1p\mathrm{}`$, an entire function $`f`$ is in $`H_R^{p,0}()`$ if, and only if, its restriction to the real line is in $`L^p(,dx)`$ and moreover $$\|f(z)\|\stackrel{~}{C}_fe^{R\|z\|}(z),$$ for some positive constant $`\stackrel{~}{C}_f`$ \[5, Theorems 6.2.4 and 6.7.1\]. This being said, for $`p=1`$ the specialization of the first part of this remark to $`k=0`$ therefore proves part of the statement in \[5, Theorem 6.8.11\], and for $`1<p<2`$ this specialization proves the first statement of \[5, Theorem 6.8.13\]. 3. The aforementioned result about entire functions shows that the specialization to $`k=0`$ of Theorem 12 is equivalent to the original Paley–Wiener theorem for the Fourier transform (see or \[18, Theorem 19.3\]). The present proof seems to be more in terms of general principles than other proofs seen in the literature.
warning/0506/cond-mat0506555.html	ar5iv	text	# Convergence and Refinement of the Wang-Landau Algorithm ## I Introduction Computational methods have been used extensively for solving complex problems in the past decades. In particular, in statistical physics equilibrium quantities of a system with many degrees of freedom are measured. The framework of statistical physics is formalized such that all equilibrium quantities can be derived from the partition function, $$𝒵(T)=\underset{\{\sigma \}}{}\text{e}^{E(\sigma )/k_\text{B}T}$$ (1) $`\sigma `$ is the state of the system, $`E`$ is the energy corresponding to $`\sigma `$, $`k_\text{B}`$ is the Boltzmann constant and $`T`$ is the temperature. The summation is over all possible states and the number of possible states is a colossal number which cannot be enumerated. Nevertheless, computational methods such as Monte Carlo techniques landau are used to sample the partition function; in particular, Metropolis importance sampling metropolis has achieved considerable success. However, new techniques are emerging and are replacing the Metropolis importance sampling especially near phase transition boundaries where the Metropolis importance sampling becomes inefficient. A class of new techniques, called the generalized ensemble methods, such as the multicanonical method berg ; lee , the broad histogram method oliveira and the flat histogram method jswang ; yamaguchi2 , were developed based on re-writing the partition function as a sum over energies $$𝒵(T)=\underset{\{\sigma \}}{}\text{e}^{E(\sigma )/k_\text{B}T}=\underset{E}{}g(E)\left[\text{e}^{E(\sigma )/k_\text{B}T}\right]$$ (2) where the partition function is reduced from a sum over all states to a sum over $``$$`N`$ energy levels. The partition function would be tractable if the energy density of states $`g(E)`$ could be calculated. Recently, a systematic, iterative, random walk method wang1 ; wang2 ; landau1 has been proposed as one of the generalized ensemble methods. Now generally known as the Wang-Landau algorithm, it has received much attention in literature and has been applied to a wide range of problems yamaguchi1 ; jain ; rathore ; yan . There have also been numerous proposed improvements and studies of the efficiency of this method schulz ; shell ; yamaguchi ; schulz1 ; zhou ; yan1 ; dayal ; trebst ; virnau ; however, there are still many unanswered questions, e.g. what determines the rate of convergence to the true density of states and is there any universality behavior related to this algorithm? In this paper, we attempt to quantify the mechanism behind the Wang-Landau method and study the effects of using different “tuning” parameters. The Wang-Landau algorithm wang1 ; wang2 ; landau1 is an iterative process in which the density of states $`g(E)`$ is modified by a factor $`f_k>1`$, and the refinement of the density of states is assured with monotonically decreasing modification factors, e.g. $`f_{k+1}=\sqrt{f_k}`$. For each time the energy level $`E`$ is visited, $`g(E)`$ is multiplied by $`f_k>1`$, and a histogram of energy is accumulated concurrently. It was proposed that the modification factor be decreased when the accumulated histogram satisfies a certain flatness condition which we call the stopping condition. In this paper we study the effects of different modification factors and stopping conditions, and derive an expression for the error term in the Wang-Landau method based on generalizations of the modification factors and stopping conditions. We shall consider arbitrary sequences of modification factors, $`f_1>\mathrm{}f_k>f_{k+1}\mathrm{}>1`$ and arbitrary sequences of corresponding stopping conditions $`\lambda _1,\lambda _2,\mathrm{}`$. The stopping conditions $`\lambda _1,\lambda _2,\mathrm{}`$ may or may not be the histogram flatness condition; other stopping conditions could be used, for example, stopping after some predetermined maximum number of Monte Carlo steps, or stopping after some number of times the random walker reaches the minimum energy state. This generalization is needed for theoretical error analysis in the next section. ## II Theoretical Error Analysis In the Wang-Landau algorithm, for each visit to an energy level $`E`$, the density of states at that energy level increases by a factor $`f_k>1`$, and the corresponding histogram increases by one. Assume that initially the unknown density of states were set to $`1`$, i.e. $`g_0(E)=1`$ $``$ $`E`$. The density of states at the end of the $`n`$th iteration is given by $$\mathrm{log}g_n(E)=\underset{k=1}{\overset{n}{}}H_k(E)\mathrm{log}(f_k)$$ (3) where $`H_k(E)`$ is the accumulated histogram and $`f_k`$ is the modification factor for the $`k`$th iteration. At this point, it is important to realize that the relative values of $`g(E)`$ are sufficient for calculating thermodynamics quantities. Hence, a constant factor can be extracted from $`g_n(E)`$, without losing any information, by a change of variable on the histograms, $`H_k(E)H_k(E)\underset{E}{\mathrm{min}}\{H_k(E)\}=\stackrel{~}{H}_k(E)`$ for $`k=1,2,\mathrm{},n`$ (4) where $`\mathrm{min}_E\{H_k(E)\}`$ is the minimum value of the accumulated histogram for the $`k`$th iteration. Then Eq. (3) becomes $$\mathrm{log}g_n(E)=\underset{k=1}{\overset{n}{}}\stackrel{~}{H}_k(E)\mathrm{log}(f_k)+\text{constant of energy}.$$ (5) The second term in Eq. (5) is independent of energy, and from here onwards we shall refer to $`\stackrel{~}{g}_n(E)`$ as the density of states without the second term in Eq. (5), i.e. $$\mathrm{log}\stackrel{~}{g}_n(E)=\underset{k=1}{\overset{n}{}}\stackrel{~}{H}_k(E)\mathrm{log}(f_k).$$ (6) To lay the foundation for deriving an expression for the error of the Wang-Landau method, we use the conjecture that the method converges to the true density of states with proper choice of parameters. ###### Conjecture 1 Let the Wang-Landau algorithm be carried out with a sequence of modification factors, $`\mathrm{}f_k>f_{k+1}>\mathrm{}f_{\mathrm{}}=1`$. There exists a sequence of stopping conditions $`\lambda _1,\lambda _2,\mathrm{}\lambda _{\mathrm{}}`$ such that, $$\underset{n\mathrm{}}{lim}\stackrel{~}{g}_n(E)=\stackrel{~}{g}_{\mathrm{}}(E)=g^{}(E)\times \text{constant}$$ (7) where $`\stackrel{~}{g}_n(E)`$ is the density of states calculated up to the $`n`$th iteration and $`g^{}(E)`$ is the true density of states. This conjecture does not give any error bounds on the density of states; it only says that, in the limit of an infinite number of iterations the Wang-Landau estimate converges to the true density of states. In addition, no constraint is imposed on the stopping conditions in the conjecture. The error term up to the $`n`$th iteration can be defined as $$\underset{E}{}\left[\mathrm{log}\stackrel{~}{g}_{\mathrm{}}(E)\mathrm{log}\stackrel{~}{g}_n(E)\right]=\underset{E}{}\underset{k=n+1}{\overset{\mathrm{}}{}}\stackrel{~}{H}_k(E)\mathrm{log}(f_k)$$ (8) An intuitive view of Eq. (8) is that, if an infinite number of iterations were performed, the exact answer would be obtained. When $`n`$ iterations were done instead, the error of $`\stackrel{~}{g}_n(E)`$ will be the sum of all the rest of the iterations that were not carried out explicitly. Define $$\mathrm{\Delta }H_k=\underset{E}{}\stackrel{~}{H}_k(E)$$ (9) where $`\mathrm{\Delta }H_k`$ is a measure of fluctuations in $`\stackrel{~}{H}_k(E)`$. By the assumption of convergence series (implied by Conjecture 7), the order of summations in Eq. (8) can be rearranged. Then, Eq. (8) becomes $$\eta _n=\underset{E}{}\left[\mathrm{log}\stackrel{~}{g}_{\mathrm{}}(E)\mathrm{log}\stackrel{~}{g}_n(E)\right]=\underset{k=n+1}{\overset{\mathrm{}}{}}\mathrm{\Delta }H_k\mathrm{log}(f_k)$$ (10) Eq. (10) shows that, assuming appropriate stopping conditions, $`\eta _n`$ depends only on the fluctuation in the histogram and the sequence of modification factors $`f_k`$. When the values of $`f_k`$ are predetermined (e.g. $`f_{k+1}=\sqrt{f_k}`$), $`\mathrm{\Delta }H_k`$ becomes the only determining factor for $`\eta _n`$. ## III Results We investigate the Monte Carlo time dependence of $`\mathrm{\Delta }H_k`$ for each iteration with the Wang-Landau method, where the subscript $`k`$ denotes the $`k`$th iteration. Simulations were performed on the ferromagnetic Ising model and on the fully frustrated Ising model with various system sizes. The Hamiltonian is $$=\underset{ij}{}J_{ij}\sigma _i\sigma _j$$ (11) where the sum is over nearest neighbors on a two dimensional square grid and $`\sigma _i`$ takes the values $`\pm 1`$. $`J_{ij}=1`$ for the ferromagnetic Ising model, and for the fully frustrated Ising model, $`J_{ij}`$ takes the value $`1`$ for every alternate horizontal nearest neighbors bonds and $`+1`$ otherwise. Fig. 1 shows the time dependence of $`\mathrm{\Delta }H_k`$ for several values of $`\mathrm{log}(f)`$; $`\mathrm{log}(f)=10^2,10^3,10^4`$ and $`10^5`$ from left to right, top to bottom respectively. We used the sequence of correction factors $`\mathrm{log}(f_{k+1})=\mathrm{log}(f_k)/1.78`$ with $`\mathrm{log}(f_1)=0.1`$, this sequence is chosen so that $`\mathrm{log}(f_{k+4})=\mathrm{log}(f_k)/10`$. These graphs were generated by performing the Wang-Landau algorithm on a $`16\times 16`$ ferromagnetic Ising model with numerical values averaged over 128 independent simulations. The Monte Carlo steps per spin, the horizontal axis of Fig. 1, are measured from the time when we decrease $`\mathrm{log}(f)`$. $`\mathrm{\Delta }H_k`$ increases initially and eventually saturates, and for smaller $`\mathrm{log}(f)`$ values, saturation values are greater and number of Monte Carlo steps required to reach saturation are larger. Because the error term given by Eq. (10) depends only on $`\mathrm{\Delta }H_k`$, any computational effort after $`\mathrm{\Delta }H_k`$ become saturated does not improve the accuracy of the final density of states $`g_n(E)`$. On the other hand, stopping the random walk before $`\mathrm{\Delta }H_k`$ becomes saturated would make the simulation less efficient because insufficient statistics are accumulated in the $`k`$th iteration and much more statistics would have to be accumulated with smaller $`\mathrm{log}(f)`$ values for subsequent iterations. An optimal algorithm is to stop the simulation as soon as $`\mathrm{\Delta }H_k`$ becomes saturated. The Wang-Landau algorithm in the original paper wang1 suggested using the histogram flatness condition as a stopping condition, but this does not guarantee optimal efficiency. It is difficult to predict the saturation value of $`\mathrm{\Delta }H_k`$ for $`k=1,\mathrm{}`$. As shown in Fig. 1, we performed a set of simulations with more Monte Carlo steps than required for $`\mathrm{\Delta }H_k`$ to reach saturation. In this way, we could measure the saturation values accurately. Fig. 2 shows a plot of saturation values versus $`\mathrm{log}(f)`$ for ferromagnetic Ising model (FMIM) and fully frustrated Ising model (FFIM). In double log scale, the data points fall on a straight line with the values of the slopes equal to $`0.491\pm 0.004`$ for $`8\times 8`$ FMIM, $`0.501\pm 0.004`$ for $`8\times 8`$ FFIM, $`0.496\pm 0.006`$ for $`16\times 16`$ FMIM and $`0.502\pm 0.008`$ for $`16\times 16`$ FFIM. To within error bars, the slopes seem to have an universal behavior $$\mathrm{max}\{\mathrm{\Delta }H_k\}\mathrm{log}(f_k)^{1/2}.$$ (12) as predicted by Zhou and Bhatt zhou . Our results suggest that the values of the slope is generic to the Wang-Landau algorithm and does not depend on system sizes and models. Certainly many more simulations on different models are needed to confirm the universality of the slope. ## IV Effects of Modification Factors We also looked at how the Wang-Landau algorithm performs with different sequences of modification factors. In the extreme case, we studied the effects of taking the limit of $`f=1`$ only after a few iterations. Assuming $`n`$ iterations were performed with large modification factors, and on the $`(n+1)`$th iteration, the modification factor is set to $`1`$. The background for implementation is as follows: Eq. (2) uses the Boltzmann weights ($`B(E,T)=\text{exp}(E/k_\text{B}T)`$) and the resulting energy distribution is, $$P(E)=g^{}(E)B(E,T)/𝒵=g^{}(E)\text{exp}(E/k_\text{B}T)/𝒵$$ (13) where $`g^{}(E)`$ is the true density of states. In general other weights can be used in summing the partition function. If one chooses $`B(E,T)=1/g_n(E)`$, then the probability distribution of $`E`$ for the $`(n+1)`$th iteration $`P_{n+1}(E)`$ will be given by, $$P_{n+1}(E)=g^{}(E)/g_n(E)𝒵$$ (14) where $`g_n(E)`$ is the density of states calculated by the $`n`$th iteration. $`𝒵`$ is is an undetermined normalization constant. The true density of states can then be estimated by the accumulated histogram of the $`(n+1)`$th iteration. $$g_{n+1}(E)=H_{n+1}(E)g_n(E)\times \text{constant}$$ (15) This is analogous to the iteration process employed in Lee’s entropic sampling lee , but we used the fast diffusion of the Wang-Landau algorithm in the early stage. Fig. 3 compares the accuracy of the Wang-Landau method with two different modification sequences for the $`32\times 32`$ ferromagnetic Ising model. The energy range was $`E/N_E[1.55,1.35]`$ where $`N_E=1024`$ is the total number of lattice sites. The vertical axis is the error of density of states defined by, $$\mathrm{\Delta }=\frac{1}{m}\underset{E}{\overset{m}{}}\left[1\frac{g_{n+1}(E)}{g^{}(E)}\right]^2$$ (16) where $`g^{}(E)`$ is the exact density of states calculated from a MATHEMATICA program provided by Beale beale . $`g_{n+1}(E)`$ is the calculated density of states and $`m`$ is the total number of energy levels in the summation over this energy range. We plot $`\mathrm{\Delta }`$ for different sequences of modification factors. Empty circles were generated with modification factors $`f_{k+1}=\sqrt{f_k}`$ (with $`f_0=\text{exp}(1)`$) and stopping when the condition $`(H_{\text{max}}H_{\text{min}})/(H_{\text{max}}+H_{\text{min}})0.1`$ is satisfied. Where $`H_{\text{max}}`$ and $`H_{\text{min}}`$ are the maximum and minimum histogram counts respectively. Filled squares were obtained with a sequence of modification factors where the limiting value of $`f=1`$ was used after 14 iterations. The arrow indicates the location which the modification factor was set to $`1`$. We measure the errors (filled squares) at several Monte Carlo steps per site after we set $`f=1`$. Error bars were obtained by averaging over several independent simulations. Accuracy increases rapidly immediately after setting $`f=1`$, but in the long run, the limiting case becomes only about twice as accurate as the other sequence. ## V Conclusion We derived an expression for the error term of the Wang-Landau algorithm. With this, we showed that the fluctuation of the accumulated histogram $`\mathrm{\Delta }H_k`$ plays a central role in the accuracy of the Wang-Landau method. We have also proposed that stopping each iteration as soon as $`\mathrm{\Delta }H_k`$ becomes saturated would be optimal. The dependence of the saturation values on the modification factor was also investigated and it was found that for the ferromagnetic Ising model and fully frustrated Ising model, $`\mathrm{max}\{\mathrm{\Delta }H_k\}\mathrm{log}(f_k)^{1/2}`$. With this equation, the saturation values of $`\mathrm{\Delta }H_k`$ for small modification factors can be predicted from values obtained with larger modification factors. Perhaps, a more efficient algorithm can be developed. We also studied the effects of using different sequences of modification factors (refinement), in which we presented the limiting case where the modification factor is set to $`1`$ after 14 iterations. There are significant improvements of efficiency for short simulations and improvements become less for longer runs. We wish to thank J. S. Wang for fruitful discussions. This work is supported by the Japan Society for Promotion of Science, the Laboratory Directed Research and Development Program of Oak Ridge National Laboratory, DOE-OS through BES-DMSE and OASCR-MICS under Contract No. DE-AC05-00OR22725 with UT-Battelle LLC, and by NSF Grant No. DMR-0341874. The computation of this work has been done using computer facilities of the Supercomputer Center, Institute of Solid State Physic, University of Tokyo (Japan) and the computer facilities of the Center for Computational Sciences, Oak Ridge National Laboratory (U.S.A.).
warning/0506/astro-ph0506432.html	ar5iv	text	# How far do they go? The outer structure of dark matter halos ## 1 Introduction More than twenty years of extensive work on cosmological N-body numerical simulations have provided numerous detailed predictions for the structure of dark matter halos in the hierarchical clustering scenario. Navarro, Frenk, & White (1997, hereafter NFW) preceded by the pioneering efforts of Quinn et al. (1986); Frenk et al. (1988); Dubinski & Carlberg (1991); Warren et al. (1992); Crone et al. (1994), suggested a simple fitting formula to describe the spherically averaged density profile of isolated dark matter halos in virial equilibrium. Since then numerous simulations were done for many relaxed halos of different masses and in different cosmologies. The NFW analytical density profile $$\rho (r)=\frac{\rho _s}{x(1+x)},xr/r_s$$ (1) has two parameters: the characteristic density $`\rho _\mathrm{s}`$ and the radius $`r_\mathrm{s}`$. Instead of these parameters, one can use the virial mass of the halo, M<sub>vir</sub>, and the concentration, $`C\text{R}\text{vir}/r_\mathrm{s}`$. Here the mass M<sub>vir</sub> and the corresponding radius R<sub>vir</sub> are defined as the mass and the radius within which the spherically averaged overdensity is equal to some specific value. For the standard cosmological model with the cosmological constant $`\mathrm{\Lambda }`$CDM and parameters $`\mathrm{\Omega }_0=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, and $`h=0.7`$ we have $`\text{M}\text{vir}=4\pi (340\rho _m)R_{\mathrm{vir}}^3/3`$, where $`\rho _m`$ is the average matter density in the Universe. For a halo with this profile, $`\rho r^1`$ as $`r0`$ and smoothly fall off as $`\rho r^3`$ at the virial radius. The concentration parameter weakly depends on the virial mass with a significant scatter comparable to the systematic change in $`C`$ over three decades in M<sub>vir</sub> (Bullock et al., 2001; Eke et al., 2001). Later simulations paid most of attention to the inner slope of the profiles. Some results favored a steeper profile than NFW density cusp with $`\rho r^{1.5}`$ (Fukushige & Makino, 1997; Moore et al., 1998; Jing & Suto, 2000; Ghigna et al., 2000). More recent simulations of halos with millions of particles within R<sub>vir</sub> seem to indicate that there is a scatter in the inner slope of the density profiles across a wide range of masses – from dwarfs to clusters. The inner slope varies between these two shapes: the NFW with an asyntotic slope of one and the steeper Moore et al. (1998) with slope 1.5 (see Klypin et al., 2001; Reed et al., 2003; Navarro et al., 2004; Diemand et al., 2004; Wambsganss et al., 2004; Tasitsiomi et al., 2004; Fukushige et al., 2004). Different approximations for density profiles were suggested and tested in the literature. Just as some other groups, we find that the 3D Sérsic three-parameter approximation gives extremely good fits for dark matter halos (Navarro et al., 2004; Merritt et al., 2005). We slightly modify this approximation by adding the average matter density of the Universe $`\rho _m`$. This term can be neglected, if one fits the density inside the virial radius. Yet, at larger distances, it gives an important contribution. The approximation can be written as $$\rho (r)=\rho _s\mathrm{exp}\left(2n[x^{1/n}1]\right)+\rho _m,xr/r_s.$$ (2) where $`n`$ is the Sérsic index. In addition to all the numerical simulations, a significant effort has been made to compare the predictions of the $`\mathrm{\Lambda }`$CDM model with the observations. This is the case of the most recent set of high quality observations of large samples of rotation curves of galaxies or the strong gravitational lensing studies which place an important upper limit on the amount of dark matter in galaxies and clusters in the inner few to tens of kiloparsec (within $`r_\mathrm{s}`$) where the need of a cuspy density profile is still subject of an exciting debate (e.g., Flores & Primack, 1994; Moore, 1994; de Blok et al., 2003; Swaters et al., 2003; Rhee et al., 2003; Keeton et al., 1998; Keeton, 2001; Broadhurst et al., 2004, and references therein). On the theoretical side, however, the origin of the shape of the dark matter halo density profile remains poorly understood. It is generally accepted that the dark matter halos are assembled by hierarchical clustering as the result of halo merging and continuous accretion. This merging scenario has motivated an interest in the analysis of the mass accretion history of the halos in conjunction with their structural properties (e.g., Wechsler et al., 2002). The systematic study of the NFW density fits to many simulated halos shows that their mass accretion history is closely correlated with the concentration parameter $`C`$ and, therefore, with the mass inside the scale radius $`r_\mathrm{s}`$ (Wechsler et al., 2002; Zhao et al., 2003; Tasitsiomi et al., 2004). These results suggest that the formation process of the dark matter halos can be generally understood by an early phase of fast mass accretion and a late phase of slow accretion of mass. In this scenario, the inner dense regions of the halos are build up early during the fast phase of mass accretion when the halo mass increases with time much faster than the expansion rate of the Universe. At later epochs, during the phase of slow mass accretion, the outer regions of the halo are built, while its inner regions stay almost intact (see Zhao et al., 2003). Despite to all this effort dedicated to the understanding of the central dense regions of the dark matter halos, very little attention has been devoted to the study of their outskirts, i.e. the regions beyond the formal virial radius. The outer parts of the halos and therefore their density profiles exhibit in these regions large fluctuations which can be understood as the result of infalling dark matter (including infalling smaller halos or substructure) or due to major mergers. In both cases the infalling material has not reached the equilibrium with the rest of the halo (see Fukushige & Makino, 2001). On the contrary, a considerable observational effort is being made to measure the mass distribution around galaxies and clusters at large distances using weak gravitational lensing (e.g., Smith et al., 2001; Guzik & Seljak, 2002; Kneib et al., 2003; Hoekstra et al., 2004; Sheldon et al., 2004, and references therein). In these cases the distances go well beyond the virial radius ranging from few hundred kpc to several Mpc. Individual field galaxies or clusters produce a small distortion of the background galaxies that allows us to measure the surface mass density profile of the dark matter (e.g., Mellier, 1999; Bartelmann & Schneider, 2001). It is customary in the weak lensing analysis to specify a dark matter halo density profile to model the projected mass profile measured with this technique. The NFW analytical formula is often adopted and extrapolated at large distances, beyond R<sub>vir</sub> with $`\rho r^3`$. This density model may not be accurate enough. The motion of satellite galaxies as a test for dark matter distribution at large radii (e.g., Zaritsky & White, 1994; Zaritsky et al., 1997; Prada et al., 2003; Brainerd, 2004; Conroy et al., 2004) is another observational method, which requires detailed theoretical predictions for outer density profiles and infall velocities. In fact, we really do not know how far the halos extend. Indeed, the NFW fitting formula was proposed and extensively tested to describe dark matter halos within R<sub>vir</sub>. This is why the NFW density fits are always done within the virial radius or even well below the virial radius ($`<0.5\text{R}\text{vir}`$), where the halos are expected to be virialized in order to avoid such non-equilibrium fluctuations. In this context, it is also surprising that the prediction of the spherical collapse model for the mean profile, which can be reasonably expected to give good predictions at sufficiently large distances ($`>2\text{R}\text{vir}`$), have not been used. In fact, the predictions of this models have only recently been worked out (see Barkana, 2004), but they have not been tested against numerical simulations. Our goal is to carry out a detailed study of the density profiles in and around collapsed objects. The paper is organized as follows. In Section 2, we give the details of the numerical simulations and the fits of the density profile of distinct galaxy-size halos. In Section 3, we have studied the shape of the density and infall velocity profiles of isolated halos up to $`23\text{R}\text{vir}`$ for different masses. We show in Section 4 that for large radii the mean density profile around dark matter halos is in excellent agreement with the predictions we have obtained via the spherical collapse model, which are somewhat different from those found by Barkana (2004). Finally, in Section 5 discussions and conclusions are given. Throughout this paper the formal virial radius is the radius within which the mean matter density is equal to 340 times the average mean matter density $`\rho _m`$ of the Universe at $`z=0`$. ## 2 Numerical simulations and density fits The simulations used in this paper are done using the Adaptive Refinement Tree (ART) code (ART, Kravtsov et al., 1997). The simulations were done for the standard $`\mathrm{\Lambda }`$CDM cosmological model with $`\mathrm{\Omega }_0=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, and $`h=0.7`$. We study halos selected from four different simulations, which parameters are listed in Table 1. The simulations cover a wide range of scales and have different mass and force resolutions. The simulation Box20 has the highest resolution, but it has only two galaxy-size halos. We use them as examples for the structure of halos simulated with very high resolution. In most of the cases we limit the analysis to well resolved halos: those should have more than $``$20,000 particles inside virial radius. The simulation Box120 has the larger volume, but its mass resolution allows us to use only halos with masses larger than $`10^{13}h^1M_{}`$. Most of the analysis of galaxy-size halos is done using simulations Box80S and Box80G. In the case of the simulation Box80G the whole $`80h^1\text{ Mpc}`$ volume was resolved with equal-mass particles. There were about 180,000 halos in the simulation. The simulation Box80S was done using particles with different masses. Only a small fraction – a $`10h^1\text{ Mpc}`$ radius region – of the box was resolved with small-mass particles. The high resolution region has an average density about equal to the mean density of the Universe. It was chosen in such a way that the halo mass function in the region is representative for the typical region of this size. For example, the region does not have a massive cluster. The most massive halo in the region has mass $`2.6\times 10^{13}h^1M_{}`$. There are 5 halos with mass larger than $`10^{13}h^1M_{}`$. Altogether, there are about 60,000 halos in this simulation. The density profile of each halo, which we study, is fit by the approximation given in eq.( 2). Each fit provides the concentration $`C`$ and the Sérsic index $`n`$. We often average density profiles of halos for some range of M<sub>vir</sub>. When doing so, we scale the radii to units of the virial radius of each halo and then average the densities. The averaged density profile is then fit again. In some cases, before we do the fitting, we also split the halo population of a given mass range into 3-4 sub-samples with a narrow range of concentrations. The parameters of the density fits together with some other properties of the halos are given in Table 2. The halos in our catalogs come from different environments. Some of them are inside the virial radii of larger halos; some have strong interactions with smaller, but still massive neighbors. We call a halo “distinct” if it does not belong to a larger halo. Most of the time we are interested in isolated halos: distinct halos, which do not have large companions. We search for halos around the given halo. If within the distance $`d\times \text{R}\text{vir}`$ the largest companion is smaller than $`\text{M}\text{vir}/m`$, then the halo is called isolated. When doing the pair-wise comparisons of the halos, we use the largest virial radius of the two halos, but we use M<sub>vir</sub> of the given halo for the test of the masses. Different isolation criteria are used. We typically use the $`d=2`$ and $`m=5`$ combination (no massive companion within 2R<sub>vir</sub>). For Milky-Way size halos with $`\text{M}\text{vir}10^{12}h^1M_{}`$ this condition typically gives 50%-60% of all distinct halos of this mass. ## 3 Halo profiles and infall velocities Figure 1 gives examples of density profiles of two halos with virial masses $`1.4\times 10^{12}h^1M_{}`$ (left panel) and $`2.6\times 10^{11}h^1M_{}`$ (right panel) in the simulation Box20. The halos have the virial radii of $`230h^1\text{ kpc}`$ and $`130h^1\text{ kpc}`$ respectively. The halos were done with very high resolution, which allows us to track the density profile below $`0.01\text{R}\text{vir}`$. The larger halo is isolated with its nearest companion being at $`3.5\text{R}\text{vir}`$. The density profile of the halo has some spikes due to substructure, but otherwise it clearly extends up to $`3\text{R}\text{vir}`$ where we see large fluctuations due to its companion. The smaller halo on the right panel has a neighbor at $`2\text{R}\text{vir}`$. So, it is not isolated. Eq.( 2) gives very good approximations for both halos. Figure 2 gives more information on the structure of these two halos. The 3D rms velocities are what one would naively expect for “normal” halos. The rms velocity first increases when we go from the center and reaches a maximum at some distance. The radius of the maximum rms velocity is smaller for the halo on the right panel. This is because it has larger concentration. At larger radii the rms velocity first declines relatively smoothly. It has some fluctuations due to substructure. At radii larger than R<sub>vir</sub> the decline stops. The average radial velocity is more interesting and to some degree is surprising. Nothing unusual inside R<sub>vir</sub>: it is practically zero with tiny ($`5\text{km s}\text{-1}`$) variations due to substructure. This is a clear sign of a virialized object. At larger distances the fluctuations in the velocity increase, but not dramatically if we compare those fluctuations with the rms velocities. The smaller halo on the right has a narrow dip ($`V_{\mathrm{rad}}20\text{km s}\text{-1}`$) at $`1.2\text{R}\text{vir}`$ apparently due to a satellite, which is moving into the halo. The surprising result is what we do not find, i.e., there is no infall on the two halos. This may seem like a fluke. Indeed, halos must grow. Their mass must increase with time. In order for the mass to increase, there should be on average negative infall velocities just outside the virial radius. Yet, we do not find those. We will later see that these two halos are not flukes, but are typical examples for halos of this mass range. What is important at this stage is that the infall velocities outside of R<sub>vir</sub> are very small. They are significantly smaller than the rms velocities. These small radial velocities indicate that halos may extend to radii significantly larger than their formal virial radii. The average density profiles of isolated halos of different masses are shown in Figure 3. The isolation criteria in this case is no massive satellite inside $`2\text{R}\text{vir}`$ ($`d=2,m=5`$; see Sec 2). The smooth density profiles, which are extremely accurately fit by eq.( 2), extend from the smallest resolved radius all the way to 2R<sub>vir</sub>. Parameters of the density fits are given in Table 2. Note that in order to reduce the range of variations along the y-axis, we plot density multiplied by $`(r/\text{R}\text{vir})^2`$. The horizontal parts of the curves in this plots correspond to density declining as $`\rho r^2`$. The density profiles are well above the average density of the Universe throughout all the radii. Even at $`5\text{R}\text{vir}`$ the average density profile is still 4-5 times larger than the mean density. The upturn at large radii tells us that the density declines less steep than $`r^2`$. The NFW approximation provides less accurate fits as shown in Figure 4. One may attribute the success of the 3D Sérsic approximation to the fact that it has three free parameters, while the NFW has only two. This is correct only to some degree. The problem with the NFW is that it has a slightly wrong shape at radii around $`r_s`$: its curvature is a bit too large. In Figure 4 this is manifested by an extended hump close to the maximum of the curves at $`r(0.050.2)\text{R}\text{vir}`$. One can shift the NFW slightly to the right and down to make the fit more accurate for most of the body of the halo ($`r>0.05\text{R}\text{vir}`$). In this case the NFW fit goes below the halo density at small radii $`r<0.05\text{R}\text{vir}`$, which sometimes was interpreted as if the central slope is steeper than -1. Overall, the 3D Sérsic approximation provides remarkable accurate fit. For $`r=(0.012)\text{R}\text{vir}`$ the errors are smaller than 5%. Figure 5 shows how the density profile depends on particular choice of the isolation criterion. In this case we selected few hundred halos in the simulation Box80G with masses $`M10^{12}h^1M_{}`$. Qualitatively the same results are found for halos with different masses. Conclusions are clear: More strict isolation conditions result in smaller density in the outer parts of halos with almost no effect inside the virial radius. Even outside of the virial radius the difference are not that large. The difference between distinct (not isolated) halos and halos, which have no massive companions inside 2R<sub>vir</sub>, are not more than a factor $`1.5`$. To large degree, this is not surprising because our isolated halos are typical, i.e., more that 1/2 of the halos in this mass range are “isolated”. One of the misconceptions, which we had before starting the analysis of the outer regions of dark matter halos is that at large distances the deviations from halo to halo are so large that it is very difficult to talk about average profile or a profile altogether. This appears to be not true. Figure 6 shows the halo to halo rms deviations from the average density profile for halos of a given mass. In order to construct the plot, we split the halo population into three ranges of concentrations and found the average and deviations for each concentration bin. Then the results of different concentrations were averaged. This splitting into concentrations is needed only for the central region $`r<0.1\text{R}\text{vir}`$ because here the average profile depends on the concentration. This plot demonstrates that there is no drastic change in the deviations at the virial radius. The deviations increase with the distance, but they are not unreasonable. Figure 7 shows the average radial velocity profile for halos of vastly different masses. We used many dozens of halos for each mass range. Just as in the case of the two individual halos in Figure 1 and 2, there is no systematic infall of material beyond formal virial radius for small virial masses. The situation is different for group- and cluster- sized halos (two top panels). For these large halos there are large infall velocities, which amplitude increases with halo mass. ## 4 Predictions from the spherical collapse model Here we obtain the predictions of the spherical collapse model (Gunn & Gott, 1972) for the $`\mathrm{𝑚𝑒𝑎𝑛}`$ halo density profile and compare them with the results found in the numerical simulations. The spherical collapse model provides a relationship between the present nonlinear enclosed density contrast, $`\delta =\rho (r)/\rho _m1`$, within a sphere of given radius $`r`$ and the enclosed linear density contrast, $`\delta _l`$, (i.e., the initial fluctuation extrapolated to the present using the linear theory) within the same sphere. This relationship along with the statistics of the initial fluctuations furnishes definite predictions for the mean density profile. We expect the spherical collapse model to give the $`\mathrm{𝑚𝑒𝑎𝑛}`$ density profile around virialized halos for sufficiently large radii - significantly larger than R<sub>vir</sub>. In fact, for $`r/\text{R}\text{vir}>2.5`$ we find a good agreement between the prediction of the spherical model and the $`\mathrm{𝑚𝑒𝑎𝑛}`$ density profiles obtained from numerical simulations for several mass ranges. By $`\mathrm{𝑚𝑒𝑎𝑛}`$ density profile we imply any representative profile for a given mass range. We consider three different representative profiles: the most probable profile, the mean profile and, what we call the typical profile. This last profile is simply the mean profile in the initial conditions spherically evolved. Our procedure to obtain the typical density profile will be presented in full detail in Betancort-Rijo et al. (2005). In essence, we use the statistics of the initial Gaussian field to obtain the mean profile of the enclosed linear density contrast, $`\delta _l(q)`$, around a proto-halo as a function of the Lagrangian distance to the center, $`q`$. We then use the spherical collapse model to obtain the present density contrast, $`\delta (r)`$, as a function of the present comoving (i.e., Eulerian) distance to the center of the halo, $`r`$. If $`Q`$ is the virial radius in Lagrangian coordinates, then the enclosed linear contrast $`\delta _l(q,Q)`$ must satisfy the condition that $`\delta _l(q=Q)=\delta _{vir},`$ where $`\delta _{vir}`$ is the linear density contrast within the virial radius R<sub>vir</sub> at the moment of virialization. The present average enclosed fractional density $`\delta (r=\text{R}\text{vir})`$ is equal to 340 (the value of overdensity $`\mathrm{\Delta }_{\mathrm{vir}}`$ used in our simulations to define the virial radius). To obtain $`\delta _{vir}`$ we must find the $`\delta _l`$ corresponding to a present density contrast equal to 340. This condition leads to a $`\delta _{vir}`$ value of 1.9. We shall later comment on this once we introduce the spherical collapse model. This condition simply ensures that the proto-halo evolves into an object which at present is virialized within the prescribed virial radius R<sub>vir</sub>. If this were the only constraint on the linear profile (equal to the initial profile except for an overall factor), it have been shown (see Section 4 in Patiri et al. 2004) that: $$\delta _l(q,Q)\delta _{vir}\mathrm{exp}\left[b\left(\left(\frac{q}{Q}\right)^21\right)\right]\delta _0(q,Q),$$ (3) where $`b`$ is a coefficient depending on R<sub>vir</sub>. It is given below. However, although this profile does not lead to wrong results, in order to achieve the accuracy required here and to obtain the correct dependence of the shape on mass of the halos we must include an additional constraint. This constraint is: $$\delta _l(q,Q)<\delta _{vir}q>Q.$$ (4) It means that there are no radii larger than R<sub>vir</sub>, where the density contrast is 340 or larger. Otherwise, the virial radius would be larger than R<sub>vir</sub>. In the large mass limit this condition becomes irrelevant, and the linear density contrast is given by eq.(3). However, in the general case the mean linear profile, $`\delta _l(q,Q)`$, is somewhat steeper than $`\delta _0(q,Q)`$: $$\delta _l(q,Q)=\delta _0(q,Q)\frac{\sigma ^{}(q,Q)\mathrm{exp}[x^2]}{1\frac{1}{2}\mathrm{erfc}[x]},$$ (5) $$x\frac{\delta _{vir}\delta _0(q,Q)}{\sqrt{2}\sigma ^{}(q,Q)}$$ (6) with $`\delta _0(q,Q)`$ given in eq.(3), and $$\sigma ^{}(q,Q)\sqrt{\sigma ^2(q)\left(\frac{\delta _l(q,Q)}{\delta _{vir}}\right)^2\sigma ^2(Q)},$$ (7) where $`\sigma (q)`$ and $`\sigma (Q)`$ are the $`rms`$ linear density fluctuations in spheres with lagrangian radii $`q`$ and $`Q`$. We use $`\sigma _8=0.9`$ as in the simulations. With these definitions we have for the previously defined function $`b(\text{R}\text{vir})`$: $$b(R_{\mathrm{vir}})=\frac{1}{4}\frac{d\mathrm{ln}\sigma (Q)}{d\mathrm{ln}Q}\|_{Q=R_{\mathrm{vir}}(340)^{1/3}}$$ (8) Using the linear density contrast eq.(5-6) we can now obtain the nonlinear density contrast, $`\delta (r)`$. For any given radius $`r`$ the nonlinear density contrast $`\delta (r)`$ is given by the solution to the equation: $$\delta _l(\delta (r))=\delta _l(q,Q),$$ (9) $$qr(1+\delta (r))^{1/3},QR_{vir}(340)^{1/3}.$$ (10) The derivation of this eq.(9) may be found in Patiri et al.(2004), although here the equation is presented in a slightly different form. The left hand side of this equation is simply the relationship, $`\delta _l(\delta )`$, between the present and the linear $`\delta `$ value in the spherical collapse model (see Sheth & Tormen 2002) evaluated at $`\delta =\delta (r)`$. For the cosmology considered here we have for $`\delta _l(\delta )`$: $`\delta _l(\delta )`$ $`=`$ $`{\displaystyle \frac{1.676}{1.68647}}[1.68647{\displaystyle \frac{1.35}{(1+\delta )^{2/3}}}`$ (11) $``$ $`{\displaystyle \frac{1.12431}{(1+\delta )^{1/2}}}+{\displaystyle \frac{0.78785}{(1+\delta )^{0.58661}}}].`$ The right hand side of eq.(9), which depends on $`\delta (r)`$ through $`q`$, is given by expression (3). Inserting eq.(5-6) into eq.(9) with $`\delta _{vir}=1.9`$, using for $`b`$ the values 0.186 and 0.254 for the two masses $`6.5\times 10^{10}h^1M_{}`$ and $`3\times 10^{12}h^1M_{}`$ discussed here, and solving for $`\delta (r)`$, we obtain the profiles for the present enclosed density contrast, $`\delta (r)`$. To obtain the density $`\delta ^{}(r)`$ at a given radius $`r`$ we simply need to take a derivative: $$(1+\delta ^{}(r))\frac{\rho (r)}{\rho _m},\delta ^{}(r)=\frac{1}{3}\frac{1}{r^2}\frac{d}{dr}r^3\delta (r),$$ (12) where $`\rho _m`$ is the mean matter density of the Universe. We must now comment on the value of $`\delta _{vir}`$ that we use. If the standard spherical collapse model were valid down to the virial radius we could use expression (11) to obtain it: $$\delta _{vir}=\delta _l(340)=1.614$$ However, we know that this is not true because bellow $`2`$ virial radius there is substantial amount of shell-crossing that render the mentioned model unappropiate. This causes $`\delta `$ to grow much more slowly with $`\delta _l`$, so that when $`\delta `$ takes the value 340, $`\delta _l`$ takes the value 1.9 (Betancort-Rijo et al., 2005). To obtain the predictions for the most probable and mean profiles, the probability distribution, $`P(\delta ,s)`$, for the value of $`\delta `$ at a given dimensionless radius $`sr/R_{vir}`$ is needed. This distribution is derived in Betancort-Rijo et al.(2005) along the same lines we have described for the the typical profile. Here we simply use it to compute the most probable, $`\delta (s)_{\mathrm{prob}}`$ profile which is is simply given by the value at which $`P(\delta ,s)`$ has its maxima. The mean profile, $`<\delta (s)>`$ is obtained from: $$<\delta (s)>=_1^{\mathrm{}}P(\delta ,s)𝑑\delta .$$ (13) In practice, since the approximation we use for $`P(\delta ,s)`$ has an unduly long tail, we have to truncate it artificially. So, we only integrate up to a $`\delta `$ value where $`P(\delta ,s)`$ has fallen to a twenty fifth of its maximum value. The results are presented in Figure 8, where we have computed, both for $`\delta `$ and for $`\delta ^{}`$, the most probable (squares) and mean (crosses) density profiles found in our simulations for two mass intervals with the mean values equal to the masses used in the theoretical derivation. We have taken $`277`$ halos in the mass range $`(6.5\pm 1.5)\times 10^{10}h^1M_{}`$ from $`Box80S`$ and $`654`$ halos in the mass range $`(3\pm 1)\times 10^{12}h^1M_{}`$ from $`Box80G`$. No isolation criteria was used. In Table 3 we list for the mean halo with mass $`<M>=3\times 10^{12}h^1M_{}`$ the estimations of the most probable and mean value of the density at different radii compare with that from the spherical collapse model for the most probable, the mean and the typical profiles. In Figure 8 one can see that beyond 2 virial radius the mean and the most probable profiles, both for $`\delta `$ and for $`\delta ^{}`$, differ considerably. This is due to the fact that for this radii the probability distribution for $`\delta `$, $`P(\delta ,s)`$, is rather wide, with a long upper tail. This can be seen in Figure 9 were this distribution is shown inside $`3.5\pm 0.05`$ virial radius for the mass $`<M>=3\times 10^{12}h^1M_{}`$. We show for comparison the theoretical prediction for $`P(\delta ,s)`$ as well as we give the most probable $`\delta _{\mathrm{max}}`$ and mean value $`<\delta >`$ of the distribution. It is apparent from Figure 8 that the $`\delta ^{}`$ profiles are steeper for smaller masses, so that they go below the background at smaller $`r/R_{vir}`$ and reach larger underdensities. The $`\delta `$ profiles are also steeper for smaller masses although the difference is, obviously, much smaller. We have found that the theoretical prediction for the typical and the most probable profile are, in general, almost indistinguishable (see Table 3). They are both found to be in very good agreement with the most probable $`\delta `$ profile found in the numerical simulations beyond two virial radii. There is also qualitative agreement between the predictions of the most probable $`\delta ^{}`$ profiles and those found in the numerical simulations. It must be noted, however, that by predicted $`\delta ^{}`$ profile we understand simply the one obtained from the corresponding $`\delta `$ profile by means of relationship given in eq.(12). Note that this is not the same as the most probable profile for $`\delta ^{}`$ (see Figure 8). This is due to the fact that the most probable $`\delta ^{}`$ value at a given $`s(r/R_{vir})`$ corresponds to a different halo than the most probable $`\delta `$ value at the same $`s`$. This explains that, while the prediction for $`\delta _{prob}`$ agrees very well with the simulations, the agreement is not so good for $`\delta _{prob}^{}`$. The $`\delta _{prob}^{}`$ obtained from $`\delta _{prob}`$ by means of expression eq.(12) is not a proper prediction but an indicative value, since we can not envisage a feasible procedure to obtain a proper prediction. On the contrary, the mean profiles $`<\delta >`$ are exactly related to $`<\delta ^{}>`$ by means of expression eq.(12). The predictions for both profiles are in good agreement with the numerical simulations, showing a much flatter profile beyond 2 virial radius than those corresponding to $`\delta _{\mathrm{prob}}`$ and $`\delta _{\mathrm{prob}}^{}`$. This agreement is remarkable given the fact that the expression used for $`P(\delta ,s)`$ is only a first approximation (see Betancort-Rijo et al. 2005). It is interesting to note, as we have previously pointed out, that larger masses have somewhat shallower profiles. In order to predict this trend correctly we must use the initial profile given by eq.(5-6). If we dropped the second constraint (that is given in eq.(4) and use in eq.(9) the initial profile given by eq.(3), which corresponds to high mass objects, the prediction would be the opposite. The reason for this being that, in this limit, the initial profile depend on mass only through $`c`$ which increases with increasing mass, thereby leading to steeper profiles for larger masses. As we stated in the introduction, the computations of the $`\mathrm{𝑚𝑒𝑎𝑛}`$ profiles (in fact, the typical profile) around halos has independently been made by Barkana (2004). He used the spherical model and, in principle, imposed on the initial profile the same constraint as we do. The computing procedure he followed was somewhat different involving some approximations. He do not give explicitly the equation defining the profile, so that accurate comparison with our results are not possible. Furthermore he used different values of $`\delta _{vir}`$, $`\mathrm{\Delta }_{vir}`$. However, his results are in good qualitative agreement with our predictions for the typical profile. We have so far considered randomly chosen halos, i.e. without isolation criteria. When halos are chosen according with an isolation criteria they differ from the randomly chosen ones in two respects. On the one hand, the probability distribution for $`\delta `$ at a given value of $`s`$ is narrower, so that most profiles cluster around the most probable one. Therefore, the difference between the mean and the most probable density profile becomes smaller. On the other hand, isolated profiles lay, on average, on somewhat more underdense environment than non-isolated ones, so that their most probable profiles are slightly steeper. Both effects may be seen by comparing Table 3 and Table 4 or the upper right pannel in Figure 8 with Figure 10 where we show the local density profile for the isolated mean halo density profile for the mass range $`3\pm 1\times 10^{12}h^1M_{}`$. In this mass range we have selected halos that do not have a companion with mass larger than 10% of the halo mass within $`4\text{R}\text{vir}`$. In total there are $`156`$ halos, i.e. one quarter of all halos in this mass range. ## 5 Conclusions and Discussions We perform a detailed study of the density profiles of isolated galaxy-size dark matter halos in high resolution cosmological simulations. We devote careful consideration mainly to the halo outer structure beyond the formal virial radius R<sub>vir</sub>. We find that the 3D Sérsic three-parameter approximation provides excellent good density fits for these dark matter halos up to 2-3R<sub>vir</sub>. These profiles do not display an abrupt change of shape beyond the virial radius. The halo-to-halo rms deviations from the average profile for halos of a given mass show that there is no a drastic change in the deviations at the virial radius. We show that these density profiles differ considerably from the NFW density profile beyond 3R<sub>vir</sub> where the density profile are slower than $`r^2`$. This result must not be seen as a contradiction when is compared with the $`r^3`$ NFW fall-off at large radii since we must remember that the NFW analytical formula was proposed and extensively tested to describe the structure of virialized dark matter halos within R<sub>vir</sub>. Although surprising, it is customary, for example, to see that the weak lensing analysis is often done with the NFW fit extrapolated to distances well beyond R<sub>vir</sub>, up to large distances of several virial radii (up to few Mpc). This approach may not be accurate enough given the results presented in this work. We also find that the isolated galaxy-size halos display all the properties of relaxed objects up to 2-3R<sub>vir</sub>. In addition to their relatively smooth density profiles seen at large radii, by studying halos average radial velocities, we find that there is no indication of systematic infall of material beyond the formal virial radius. The dark matter halos in this mass range do not grow as one naively may expect through a steady accretion of satellites, i.e., on average there is no mass infall. This is strikingly different for more massive halos, such as group- and cluster-sized halos which exhibit large infall velocities outside of the formal virial radius. For large halos the amplitude of the infall velocities increases with halo mass. For larger radii beyond 2-3 formal virial radius we combine the statistics of the initial fluctuations with the spherical collapse model to obtain predictions of the mean halo density profiles for halos with different masses. We consider two possibilities: the most probable and the mean density profiles. We find that the most probable profile obtained from our simulations is in excellent agreement with the predictions from the spherical collapse model beyond 2-3 virial radius. For the mean density profile the predictions are not so accurate. This is due to the fact that the approximation, which we are using for the distribution of $`\delta `$ at a given radius, has an artificially long tail (we are presently working on a better approximation). Even so, the predictions are qualitatively good and quantitatively quite acceptable. We think that the discrepancies between the data and the predictions at radii smaller than 2-3 virial radii are due to the fact that these inner shells are affected by the shell-crossing. We hope that an appropiate treatment of this circunstance will lead to accurate predictions for all radii, so that the mean spherically averaged profiles may be understood in terms of the spherical collapse. We find the results presented here very encouraging in this respect and we are currently working along this line. We thank Yehuda Hoffman and Simon White for stimulating discussions. F.P. is a Ramón y Cajal Fellow at the IAA (CISC). F.P. want to thank the hospitality and financial support of the Instituto de Astrofisica de Canarias and the New Mexico State University where part of this work was done. F.P., S.P. and S.G. would like to thank Acciones Integradas for supporting the German-Spanish collaboration. S.G. acknowledges support by DAAD. A.K. acknowledges support of NASA and NSF grants to NMSU. Computer simulations have been done at the LRZ Munich, NIC Jülich and the NASA Ames. Part of the data analysis have been carry out at CESGA.
warning/0506/hep-ph0506074.html	ar5iv	text	# Test of Lorentz Invariance in Electrodynamics Using Rotating Cryogenic Sapphire Microwave Oscillators ## Abstract We present the first results from a rotating Michelson-Morley experiment that uses two orthogonally orientated cryogenic sapphire resonator-oscillators operating in whispering gallery modes near 10 GHz. The experiment is used to test for violations of Lorentz Invariance in the frame-work of the photon sector of the Standard Model Extension (SME), as well as the isotropy term of the Robertson-Mansouri-Sexl (RMS) framework. In the SME we set a new bound on the previously unmeasured $`\stackrel{~}{\kappa }_e^{ZZ}`$ component of $`2.1(5.7)\times 10^{14}`$, and set more stringent bounds by up to a factor of 7 on seven other components. In the RMS a more stringent bound of $`0.9(2.0)\times 10^{10}`$ on the isotropy parameter, $`P_{MM}=\delta \beta +\frac{1}{2}`$ is set, which is more than a factor of 7 improvement. The Einstein Equivalence Principle (EEP) is a founding principle of relativity Will . One of the constituent elements of EEP is Local Lorentz Invariance (LLI), which postulates that the outcome of a local experiment is independent of the velocity and orientation of the apparatus. The central importance of this postulate has motivated tremendous work to experimentally test LLI. Also, a number of unification theories suggest a violation of LLI at some level. However, to test for violations it is necessary to have an alternative theory to allow interpretation of experiments Will , and many have been developed Robertson ; MaS ; LightLee ; Ni ; Kosto1 ; KM . The kinematical Roberson-Mansouri-Sexl (RMS) Robertson ; MaS framework postulates a simple parameterization of the Lorentz transformations with experiments setting limits on the deviation of those parameters from their values in special relativity (SR). Because of their simplicity they have been widely used to interpret many experiments Brillet ; Wolf ; Muller ; WolfGRG . More recently, a general Lorentz violating extension of the standard model of particle physics (SME) has been developed Kosto1 whose Lagrangian includes all parameterized Lorentz violating terms that can be formed from known fields. This work presents first results of a rotating lab experiment using cryogenic microwave oscillators. Previous non-rotating experiments Lipa ; Muller ; Wolf04 relied on the earth’s rotation to modulate a Lorentz violating effect. This is not optimal for two reasons. Firstly, the sensitivity is proportional to the noise of the oscillators at the modulation frequency, typically best for periods between 10 and 100 seconds. Secondly, the sensitivity is proportional to the square root of the number of periods of the modulation signal, therefore taking a relatively long time to acquire sufficient data. Thus, by rotating the experiment the data integration rate is increased and the relevant signals are translated to the optimal operating regime Mike . Our experiment consists of two cylindrical sapphire resonators of 3 cm diameter and height supported by spindles within superconducting niobium cavities Giles , and are oriented with their cylindrical axes orthogonal to each other in the horizontal plane. Whispering gallery modes wgmode are excited near 10 GHz, with a difference frequency of 226 kHz. The frequencies are stabilized using Pound locking, and amplitude variations are suppressed using an additional control circuit. A detailed description of such oscillators can be found in Mann ; Hartnett . The resonators are mounted in a common copper block, which provides common mode rejection of temperature fluctuations. The structure is in turn mounted inside two successive stainless steel vacuum cylinders from a copper post, which provides the thermal connection between the cavities and the liquid helium bath. A foil heater and carbon-glass temperature sensor attached to the copper post controls the temperature set point to 6 K with mK stability. A schematic of the rotation system is shown in Fig.1. The cryogenic dewar along with the room temperature oscillator and control electronics is suspended within a ring bearing. A multiple ”V” shaped suspension made from elastic cord avoids high Q-factor pendulum modes by ensuring that the cord has to stretch and shrink (providing damping) for horizontal and vertical motion. The rotation system is driven by a microprocessor controlled stepper motor. A commercial 18 conductor slip ring connector, with a hollow through bore, transfers power and various signals to and from the rotating experiment. A mercury based rotating coaxial connector transmits the difference frequency to a stationary frequency counter referenced to an Oscilloquartz oscillator. The data acquisition system logs the difference frequency as a function of orientation, as well as monitoring systematic effects including the temperature of the resonators, liquid helium bath level, ambient room temperature, oscillator control signals, tilt, and helium return line pressure. Inside the sapphire crystals standing waves are set up with the dominant electric and magnetic fields in the axial and radial directions respectively, corresponding to a Poynting vector around the circumference. The experimental observable is the difference frequency, and to test for Lorentz violations the perturbation of the observable with respect to an alternative test theory must be derived. For example, in the photon sector of the SME this may be calculated to first order as the integral over the non-perturbed fields (Eq. (34) of KM ), and expressed in terms of 19 independent variables (discussed in more detail later). The change in orientation of the fields due to the lab rotation and Earth’s orbital and sidereal motion induces a time varying modulation of the difference frequency, which is searched for in the experiment. Alternatively, with respect to the RMS framework, we analyze the change in resonator frequency as a function of the Poynting vector direction with respect to the velocity of the lab through the cosmic microwave background (CMB). The RMS parameterizes a possible Lorentz violation by a deviation of the parameters ($`\alpha ,\beta ,\delta `$) from their SR values ($`\frac{1}{2},\frac{1}{2},0`$). Thus, a complete verification of LLI in the RMS framework Robertson ; MaS requires a test of (i) the isotropy of the speed of light ($`P_{MM}=\delta \beta +\frac{1}{2}`$), a Michelson-Morley (MM) experiment MM , (ii) the boost dependence of the speed of light ($`P_{KT}=\beta \alpha 1`$), a Kennedy-Thorndike (KT) experiment KT and (iii) the time dilation parameter ($`P_{IS}=\alpha +\frac{1}{2}`$), an Ives-Stillwell (IS) experiment IS ; Saat . Because our experiment compares two cavities it is only sensitive to $`P_{MM}`$. Fig.2 shows typical fractional frequency instability of the 226 kHz difference with respect to 10 GHz, and compares the instability when rotating and stationary. A minimum of $`1.6\times 10^{14}`$ is recorded at 40s. Rotation induced systematic effects degrade the stability up to 18s due to signals at the rotation frequency of $`0.056Hz`$ and its harmonics. We have determined that tilt variations dominate the systematic by measuring the magnitude of the fractional frequency dependence on tilt and the variation in tilt at twice the rotation frequency, $`2\omega _R(0.11Hz)`$, as the experiment rotates. We minimize the effect of tilt by manually setting the rotation bearing until our tilt sensor reads a minimum at $`2\omega _R`$. The latter data sets were up to an order of magnitude reduced in amplitude as we became more experienced at this process. The remaining systematic signal is due to the residual tilt variations, which could be further annulled with an automatic tilt control system. It is still possible to be sensitive to Lorentz violations in the presence of these systematics by measuring the sidereal, $`\omega _{}`$ and semi-sidereal, $`2\omega _{}`$ sidebands about $`2\omega _R`$, as was done in Brillet . The amplitude and phase of a Lorentz violating signal is determined by fitting the parameters of Eq.1 to the data, with the phase of the fit adjusted according to the test theory used. $$\frac{\mathrm{\Delta }\nu _0}{\nu _0}=A+Bt+\underset{i}{}C_i\mathrm{cos}(\omega _it+\phi _i)+S_i\mathrm{sin}(\omega _it+\phi _i)$$ (1) Here $`\nu _0`$ is the average unperturbed frequency of the two sapphire resonators, and $`\mathrm{\Delta }\nu _0`$ is the perturbation of the 226 kHz difference frequency, $`A`$ and $`B`$ determine the frequency offset and drift, and $`C_i`$ and $`S_i`$ are the amplitudes of a cosine and sine at frequency $`\omega _i`$ respectively. In the final analysis we fit 5 frequencies to the data, $`\omega _i=(2\omega _R,2\omega _R\pm \omega _{},2\omega _R\pm 2\omega _{})`$, as well as the frequency offset and drift. The correlation coefficients between the fitted parameters are all between $`10^2`$ to $`10^5`$. Since the residuals exhibit a significantly non-white behavior, the optimal regression method is weighted least squares (WLS) Wolf04 . WLS involves pre-multiplying both the experimental data and the model matrix by a whitening matrix determined by the noise type of the residuals of an ordinary least squares analysis. We have acquired 5 sets of data over a period of 3 months beginning December 2004, totaling 18 days. The length of the sets (in days) and size of the systematic are ($`3.6,2.3\times 10^{14}`$), ($`2.4,2.1\times 10^{14}`$), ($`1.9,2.6\times 10^{14}`$), ($`4.7,1.4\times 10^{15}`$), and ($`6.1,8.8\times 10^{15}`$) respectively. We have observed leakage of the systematic into the neighboring side bands due to aliasing when the data set is not long enough or the systematic is too large. Fig.3 shows the total amplitude resulting from a WLS fit to 2 of the data sets over a range of frequencies about $`2\omega _R`$. It is evident that the systematic of data set 1 at $`2\omega _R`$ is affecting the fitted amplitude of the sidereal sidebands $`2\omega _R\pm \omega _{}`$ due to its relatively short length and large systematics. By analyzing all five data sets simultaneously using WLS the effective length of the data is increased, reducing the width of the systematic sufficiently as to not contribute significantly to the sidereal and semi-sidereal sidebands. In the photon sector of the SME KM , Lorentz violating terms are parameterized by 19 independent components, which are in general grouped into three traceless and symmetric $`3\times 3`$ matrices ($`\stackrel{~}{\kappa }_{e+}`$, $`\stackrel{~}{\kappa }_o`$, and $`\stackrel{~}{\kappa }_e`$), one antisymmetric matrix($`\stackrel{~}{\kappa }_{o+}`$) and one additional scalar, which all vanish when LLI is satisfied. To derive the expected signal in the SME we use the method of KM ; WolfGRG to calculate the frequency of each resonator in the SME and in the resonator frame. We then transform to the standard celestial frame used in the SME KM taking into account the rotation in the laboratory frame in a similar way to TobarPRD . The resulting relation between the parameters of the SME and the $`C_i`$ and $`S_i`$ coefficients are given in Tab.1 which, for short data sets, were calculated using the leading order expansion at the annual phase position of the data. The 10 independent components of $`\stackrel{~}{\kappa }_{e+}`$ and $`\stackrel{~}{\kappa }_o`$ have been constrained by astronomical measurements to $`<2\times 10^{32}`$ KM ; Kost01 . Seven components of $`\stackrel{~}{\kappa }_e`$ and $`\stackrel{~}{\kappa }_{o+}`$ have been constrained in optical and microwave cavity experiments Muller ; Wolf04 at the $`10^{15}`$ and $`10^{11}`$ level respectively, while the scalar $`\stackrel{~}{\kappa }_{tr}`$ component recently had an upper limit set of $`<10^4`$ TobarPRD . The remaining $`\stackrel{~}{\kappa }_e^{ZZ}`$ component could not be previously constrained in non-rotating experiments Muller ; Wolf04 . In contrast, our rotating experiment is sensitive to $`\stackrel{~}{\kappa }_e^{ZZ}`$. However, it appears only at $`2\omega _R`$, which is dominated by systematic effects. From our combined analysis of all data sets, and using the relation to $`\stackrel{~}{\kappa }_e^{ZZ}`$ given in Tab.1, we determine a value for $`\stackrel{~}{\kappa }_e^{ZZ}`$ of $`4.1(0.5)\times 10^{15}`$. However, since we do not know if the systematic has canceled a Lorentz violating signal at $`2\omega _R`$, we cannot reasonably claim this as an upper limit. Since we have five individual data sets, a limit can be set by treating the $`C_{2\omega _R}`$ coefficient as a statistic. The phase of the systematic depends on the initial experimental conditions, and is random across the data sets. Thus, we have five values of $`C_{2\omega _R}`$, ($`\{4.2,11.4,21.4,1.3,8.1\}`$ in $`10^{15}`$). If we take the mean of these coefficients, the systematic signal will cancel if its phase is random, but the possible Lorentz violating signal (with constant phase) will not. Thus a limit can be set by taking the mean and standard deviation of the five coefficient of $`C_{2\omega _R}`$. This gives a more conservative bound of $`2.1(5.7)\times 10^{14}`$, which includes zero. Our experiment is also sensitive to all other seven components of $`\stackrel{~}{\kappa }_e`$ and $`\stackrel{~}{\kappa }_{o+}`$ (see Tab.1) and improves present limits by up to a factor of 7, as shown in Tab.2. In the RMS frame-work, a frequency shift due to a putative Lorentz violation is given by Eq.2 Wolf ; WolfGRG , $$\frac{\mathrm{\Delta }\nu _0}{\nu _0}=\frac{P_{MM}}{2\pi c^2}[(𝐯.\widehat{\theta }_\mathrm{𝟏})^2d\phi _1(𝐯.\widehat{\theta }_\mathrm{𝟐})^2d\phi _2]$$ (2) Where $`𝐯`$ is the velocity of the preferred frame wrt the CMB, $`\widehat{\theta }_j`$ is the unit vector in the direction of the azimuthal angle (direction of propagation) of each resonator (labeled by subscripts 1 and 2), and $`\phi `$ is the azimuthal variable of integration in the cylindrical coordinates of each resonator. Perturbations due to Lorentz violations occur at the same five frequencies as the SME, but for the RMS analysis we do not consider the $`2\omega _R`$ frequency due to the large systematic, as we only need to put a limit on one parameter. The dominant coefficients are due to only the cosine terms with respect to the CMB right ascension, $`Cu_i`$, which are shown in Tab.3. In conclusion, we set bounds on 7 components of the SME photon sector (Tab.2) and $`P_{MM}`$ (Tab.3) of the RMS framework, which are up to a factor of 7 more stringent than those obtained from previous experiments. We have also set an upper limit \[$`2.1(5.7)\times 10^{14}`$\] on the previously unmeasured SME component $`\stackrel{~}{\kappa }_e^{ZZ}`$. To further improve these results, tilt and environmental controls will be implemented to reduce systematic effects. To remove the assumption that $`\stackrel{~}{\kappa }_{o+}`$ and $`\stackrel{~}{\kappa }_e`$ do not cancel each other, data integration will continue for more than a year. Note added: Two other concurrent experiments have also set some similar limits Ant ; Herr . ###### Acknowledgements. This work was funded by the Australian Research Council.
warning/0506/gr-qc0506075.html	ar5iv	text	# General Relativity in the Undergraduate Physics Curriculum ## I Introduction Einstein’s 1915 relativistic theory of gravity — general relativity — will soon be a century old. It is the classical theory of one of the four fundamental forces. It underlies our contemporary understanding of the big bang, black holes, pulsars, quasars, X-ray sources, the final destiny of stars, gravitational waves, and the evolution of the universe itself. It is the intellectual origin of many of the ideas at play in the quest for a unified theory of the fundamental forces that includes gravity. The heart of general relativity is one of the most beautiful and revolutionary ideas in modern science — the idea that gravity is the geometry of curved four-dimensional spacetime. General relativity and quantum mechanics are usually regarded as the two greatest developments of twentieth-century physics. Yet, paradoxically, general relativity — so well established, so important for several branches of physics, and so simple in its basic conception — is often not represented anywhere in the undergraduate physics curriculum. An informal survey by William Hiscock His05 of the course offerings of 32 mid-western research universities found only a handful that offered an intermediate (junior/senior) course in general relativity as part of the undergraduate physics curriculum. This has the consequence that many students see gravity first in the context of planetary orbits in basic mechanics and next, if at all, in an advanced graduate course designed in part for prospective specialists in the subject. There might have been an argument for such an organization half a century ago. But there is none today in an era when gravitational physics is increasingly important, increasingly topical, increasingly integrated with other areas of physics, and increasingly connected with experiment and observation. In the author’s opinion, every undergraduate physics major should have an opportunity to be introduced to general relativity. Its importance in contemporary physics is not the only reason for introducing undergraduates to general relativity. There are others: First, the subject excites interest in students. Warped spacetime, black holes, and the big bang are the focus both of contemporary research and of popular scientific fascination. Students specializing in physics naturally want to know more. A further argument for undergraduate general relativity is accessibility. As I hope to show in this paper, a number of important phenomena of gravitational physics can be efficiently introduced with just a basic background in mechanics and a minimum of mathematics beyond the usual advanced calculus tool kit. Other subjects of great contemporary importance such as high temperature superconductivity or gauge theories of the strong interactions require much more prerequisite information. General relativity can be made accessible to both students and faculty alike at the undergraduate level. It is probably fruitless to speculate on why a subject as basic, accessible, and important as general relativity is not taught more widely as part of the undergraduate physics curriculum. Limited time, limited resources, inertia, tradition, and misconceptions may all play a role. Certainly it is not a lack of textbooks. Refs. Ber93 Tou97 are a partial list of texts known to the author<sup>1</sup><sup>1</sup>1This list consists of texts known to the author, published after 1975, and judged to be introductory. It does not pretend to be either complete or selective, nor is it a representation that the texts are readily available. that treat general relativity in some way at an introductory level. Available time is one of the obstacles to introducing general relativity at the undergraduate level. The deductive approach to teaching this subject (as for most others) is to assemble the necessary mathematical tools, motivate the field equations, solve the equations in interesting circumstances, and compare the predictions with observation and experiment. This ‘math first’ order takes time to develop for general relativity which may not be available to either students or faculty. This article describes a different, ‘physics first’ approach to introducing general relativity at the junior/senior level. Briefly, the simplest physically relevant solutions to the Einstein equation are introduced first, without derivation, as curved spacetimes whose properties and observable consequences can be explored by a study of the motion of test particles and light rays. This brings the student to interesting physical phenomena as quickly as possible. It is the part of the subject most directly connected to classical mechanics and the part that requires a minimum of ‘new’ mathematical ideas. Later the Einstein equation can be motivated and solved to show where the solutions come from. When time is limited this is a surer and more direct route to getting at the applications of general relativity that are important in contemporary science. Section II expands very briefly on the importance of general relativity in contemporary physics. Section III outlines the basic structure of the subject. Sections IV and V describe the ‘math first’ and ‘physics first’ approaches to introducing general relativity to undergraduate physics majors. This is not an even-handed comparison. The ‘math first’ approach is described only to contrast it with the ‘physics first’ approach which is advocated in this paper. Section VI illustrates how ideas from classical mechanics can be used to calculate important effects in general relativity. Section VII reports the personal experiences of the author in using the ‘physics first’ method. ## II Where is General Relativity Important? Gravity is the weakest of the four fundamental forces at accessible energy scales. The ratio of the gravitational force to the electric force between two protons separated by a distance $`r`$ is (in Gaussian electromagnetic units) $$\frac{F_{\mathrm{grav}}}{F_{\mathrm{elec}}}=\frac{Gm_p^2/r^2}{e^2/r^2}=\frac{Gm_p^2}{e^2}10^{40}.$$ (2.1) Gravity might thus seem to be negligible. But three other facts explain why it is important and where it is important. First, gravity is a universal force coupling to all forms of mass and energy. Second, gravity is a long-range force in contrast to the weak and strong forces which are characterized by nucleus-size ranges and below. Third, and most importantly, gravity is unscreened. There is no negative “gravitational charge”; mass is always positive. These three facts explain why gravity is the dominant force governing the structure of the universe on the largest scales of space and time — the scales of astrophysics and cosmology. The strong and weak forces are short range. The relatively much greater strength of electromagnetic forces ensures that charges will be screened in an electrically neutral universe like ours. Only gravity is left to operate on very large scales. Relativistic gravity — general relativity — is important for an object of mass $`M`$ and size $`R`$ when $$q\frac{GM}{Rc^2}1.$$ (2.2) Neutron stars $`(q.1)`$ and black holes $`(q.5)`$ are relativistic objects by this rough criterion. So is our universe $`(q1)`$ if we take $`R`$ to be the present Hubble distance and $`M`$ to be the mass within it. Figure 1 displays some phenomena for which relativity is important and ones for which it is not. General relativity can sometimes be important even when $`q`$ is small provided compensating observational precision can be achieved. For the Sun $`q10^6`$, yet the solar system is the domain of the precision tests that confirm general relativity to as much as 1 part in $`10^5`$ Cassini . For the Earth $`q10^9`$, yet general relativistic effects are important for the operation of the Global Positioning System (GPS) Ash02 . Relativistic gravity is also important on the smallest scales considered in contemporary physics — those of quantum gravity. These are characterized by the Planck length $`\mathrm{}`$ $$\mathrm{}_{\mathrm{Pl}}=(G\mathrm{}/c^3)^{\frac{1}{2}}10^{33}\mathrm{cm}.$$ (2.3) This is much, much smaller than even the scale of the strong interactions $`10^{13}`$ cm. Yet, this is the scale which many contemporary explorers believe will naturally characterize the final theory unifying the four fundamental forces including gravity. This is the characteristic scale of string theory. This is the scale that will characterize the union of the two great developments of twentieth century physics — general relativity and quantum mechanics. The important point for this discussion is that the last few decades have seen dramatic growth in observational data on the frontier of the very large, and an equally dramatic growth in theoretical confidence in exploring the frontier of the very small. Black holes, for example, are no longer a theorist’s dream. They have been identified at the center of galaxies (including our own) and in X-ray binaries. They are central to the explanations of the most energetic phenomena in the universe such as active galactic nuclei. On even larger scales, it is now a commonplace observation that cosmology has become a data driven science. Cosmological parameters once uncertain by orders of magnitude have been determined to accuracies of 10% Key ; Wmap . Adventures into Planck scale physics may be mostly in the minds of theorists, but the quest for a unified theory of the fundamental forces including gravity is being pursued with impressive vigor and confidence by a large community. Indeed, at the big bang where large and small are one, we should eventually see direct evidence of Planck scale physics. For these reasons general relativity is increasingly central to today’s physics. It is increasingly topical, increasingly connected with experiment and observation, and increasingly integrated with other branches of physics<sup>2</sup><sup>2</sup>2For more on the importance of general relativity in contemporary physics, see e.g. cpg99 . ## III Key Ideas in General Relativity This section sketches a few key ideas in general relativity for those who may not be familiar with the theory. The intent is not to offer an exposition of these ideas. That, after all, is properly the task of texts on the subject. Rather, the purpose is merely to mention ideas that will occur in the subsequent discussion and to illustrate the simplicity of the conceptual structure of the subject. It’s potentially misleading to summarize any subject in physics in terms of slogans. However, the following three roughly stated ideas are central to general relativity. * Gravity is Geometry. Phenomena familiarly seen as arising from gravitational forces in a Newtonian context are more generally due to the curvature of geometry of four-dimensional spacetime. * Mass-Energy is the Source of Spacetime Curvature. Mass is the source of spacetime curvature and, since general relativity incorporates special relativity, any form of energy is also a source of spacetime curvature. * Free Mass Moves on Straight Paths in Curved Spacetime. In general relativity, the Earth moves around the Sun in the orbit it does, not because of a gravitational force exerted by the Sun, but because it is following a straight path in the curved spacetime produced by the Sun. Making these ideas more precise and more explicit is an objective of any course in general relativity. We mention a few steps toward this objective here. Points in four-dimensional spacetime can be located by four coordinates $`x^\alpha ,\alpha =0,1,2,3`$. Coordinates are arbitrary provided they label points uniquely. Generally, several different coordinate patches are required to label all the points in spacetime. The geometry of a spacetime is specified by giving the metric, $`g_{\alpha \beta }(x)`$, where the $`(x)`$ indicates that the metric is generally a function of all four coordinates. The metric determines the squared distance $`ds^2`$ in four-dimensional spacetime between points separated by infinitesimal coordinate intervals $`dx^\alpha `$. Specifically, $$ds^2=g_{\alpha \beta }(x)dx^\alpha dx^\beta $$ (3.1) where a double sum over $`\alpha `$ and $`\beta `$ from 0 to 3 is implied. Integration of the $`ds`$ specfied by this expression gives the distance along curves. Metrics satisfy the Einstein equation $$R_{\alpha \beta }\frac{1}{2}g_{\alpha \beta }R=\frac{8\pi G}{c^4}T_{\alpha \beta }$$ (3.2) relating a measure of curvature on the left hand side to the energy-momentum tensor of matter on the right. This Einstein equation comprises 10 non-linear, partial differential equations for the metric $`g_{\alpha \beta }(x)`$. An important example of a solution to the Einstein equation is the Schwarzschild geometry giving the metric in the empty space outside a spherically symmetric black hole or star. In standard Schwarzschild spherical coordinates $`x^\alpha =(t,r,\theta ,\varphi )`$ this is $`ds^2=`$ $`\left(1{\displaystyle \frac{2GM}{c^2r}}\right)(cdt)^2+\left(1{\displaystyle \frac{2GM}{c^2r}}\right)^1dr^2`$ $`+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)`$ (3.3) where $`M`$ is the mass of the black hole or star. This, to an excellent approximation, describes the curved spacetime outside our Sun. Test particles with masses too small to affect the ambient geometry move on straight paths in it. More precisely, they move between any two points, $`A`$ and $`B`$, in spacetime on a world line (curve) of stationary proper time $`\tau `$. The proper time along a world line is the distance along it measured in time units. Thus, $`d\tau ^2=ds^2/c^2`$. (The negative sign is so $`d\tau ^2`$ is positive along the world lines of particles which always move with less than the speed of light). Curves of stationary proper time are called geodesics. The world line of a particle through spacetime from $`A`$ to $`B`$ can be described by giving its coordinates $`x^\alpha (\lambda )`$ as a function of any parameter that takes fixed values on the end points. For instance using the Schwarzschild metric (3.3), the principle of stationary proper time takes the form: $`\delta {\displaystyle _A^B}`$ $`d\tau =\delta {\displaystyle _A^B}d\lambda {\displaystyle \frac{1}{c^2}}[(1{\displaystyle \frac{2GM}{c^2r}})(c\dot{t})^2`$ $`(1{\displaystyle \frac{2GM}{c^2r}})^1\dot{r}^2r^2(\dot{\theta }^2+\mathrm{sin}^2\theta \dot{\varphi }^2)]^{\frac{1}{2}}=0`$ (3.4) where a dot denotes a derivative with respect to $`\lambda `$ and $`\delta `$ means the first variation as in classical mechanics. The variational principle (3.4) for stationary proper time has the same form as the variational principle for stationary action in classical mechanics. The Lagrangian $`L(\dot{x}^\alpha ,x^\alpha )`$ is the integrand of (3.4). Lagrange’s equations are the geodesic equations of motion. Their form can be made especially simple by choosing proper time $`\tau `$ for the parameter $`\lambda `$. From them one can deduce the conservation of energy, the conservation of angular momentum, and an effective equation for radial motion in the Schwarzschild geometry as we will describe in Section VI \[cf.(6.1)\]. With that, one can calculate a spectrum of phenomena ranging from the precession of planetary orbits to the collapse to a black hole. When extended to light rays and implemented in appropriate metrics, the geodesic equations are enough to explore most of the important applications of relativistic gravity displayed in Table 1. ## IV Teaching General Relativity — Math First The deductive approach to teaching many subjects in physics is to 1. Introduce the necessary mathematical tools; 2. Motivate and explain the basic field equations; 3. Solve the field equations in interesting circumstances; 4. Apply the solutions to make predictions and compare with observation and experiment. For electromagnetism, (1)–(4) are, e.g. (1) Vector calculus, (2) Maxwell’s equations, (3) boundary value problems, the fields of point particles, radiation fields, etc., (4) charged particle motion, circuits, wave guides and cavities, antennas, dielectric and magnetic materials, magnetohydrodynamics — a list that could very easily be extended. For gravitation (1)–(4) are e.g. (1) differential geometry, (2) the Einstein equation and the geodesic equation, (3) the solutions for spherical symmetry, cosmological models, gravitational waves, relativistic stars, etc. Table 1 lists some of the applications of general relativity that constitute (4). This deductive order of presentation is logical; it is the order used by the great classic texts LL62 ; MTW70 ; Wei72 ; Wal84 ; and it is the order used in standard graduate courses introducing the subject at an advanced level. But the deductive order does have some drawbacks for an elementary introduction to physics majors in a limited time. Differential geometry is a deep and beautiful mathematical subject. However, even an elementary introduction to its basic ideas and methods are not a part of the typical advanced calculus tool kit acquired by physics majors in their first few years. This is ‘new math’. In contrast, the vector calculus central to electromagnetism is part of this tool kit. It is possible at the undergraduate level to give an introduction to the basic mathematical ideas of manifolds, vectors, dual vectors, tensors, metric, covariant derivative, and curvature. Indeed, many students feel empowered by learning new mathematics. But it does take time. It also must be practiced. The author’s experience is that many students at this level need considerable exercise before they are able to accurately and efficiently manipulate tensorial expressions and feel at home with the four-dimensional mathematical concepts necessary to formulate Einstein’s equation. When time is limited, pursuing the deductive order may leave little available for the interesting applications of general relativity. Further, solving the Einstein equation to exhibit physically relevant spacetime geometries is a difficult matter. Their non-linear nature means that there is no known general solution outside of linearized gravity. Each new situation, e.g. spherical symmetry, homogeneous and isotropic cosmological models, gravitational plane waves, is typically a new problem in applied mathematics. Deriving the solutions only adds to the time expended before interesting applications can be discussed. Many of the successful introductory texts in general relativity follow this logical order at various levels of compromise. In the author’s opinion, an outstanding example is Bernard Schutz’s classic A First Course in General Relativity Sch85 . In the next section we consider a different way of introducing general relativity to undergraduates. ## V Teaching General Relativity — Physics First Electricity and magnetism are not usually presented in introductory (freshman) courses in the deductive order described in the previous section. Specifically, we do not usually first develop vector calculus, then exhibit Maxwell’s equations, then solve for the fields of charges, currents, and radiation, and finally apply these to realistic electromagnetic phenomena. Rather, the typical course posits the fields of the simplest physically relevant examples, for instance the electric field of a point charge, the magnetic field of a straight wire, and the electromagnetic plane wave. These are used to build understanding of fields and their interaction with charges for immediate application to demonstrable electromagnetic phenomena. Maxwell’s ten partial differential equations and their associated gauge and Lorentz invariances are better appreciated later, usually in a more advanced course. General relativity can be efficiently introduced at an intermediate (junior/senior) level following the same ‘physics first’ model used in introducing electromagnetism. Specifically: 1. Exhibit the simplest physically important spacetime geometries first, without derivation; 2. Derive the predictions of these geometries for observation by a study of the orbits of test particles and light rays moving in them; 3. Apply these predictions to realistic astrophysical situations and compare with experiment and observation; 4. Motivate the Einstein equation and solve it to show where the spacetime geometries posited in (1) come from. These are essentially the same four elements that comprise the deductive approach described in the previous section, but in a different order. That order has considerable advantages for introducing general relativity at an intermediate level as we now describe: ### V.1 Indications Less ‘New Math’ Up Front: To exhibit a spacetime geometry, the only ‘new math’ ideas required are the metric and its relation to distances in space and time. To analyze the motion of test particles in these geometries, only the notions of four-vectors and geodesics are needed. These three new mathematical ideas are enough to explain in detail a wide range of physical phenomena, such as most of those in Table 1. Further, these three new mathematical ideas are among the simplest parts of a relativist’s tool kit to introduce at an intermediate level. Four-vectors are often familiar from special relativity. Geodesics viewed as curves of extremal proper time are special cases of Lagrangian mechanics. The idea of a metric can be motivated from the theory of surfaces in three-dimensional flat space. A general theory of tensors as linear maps from vectors into the real numbers is not required because only one tensor — the metric — is ever used. The Simplest Spacetimes are the Most Physically Relevant: * The Sun is approximately spherical. * Spherical black holes exhibit many characteristic properties of the most general black hole. * The universe is approximately homogeneous and isotropic on scales above 100 megaparsecs. * Detectable gravitational waves are weak. These four facts mean that the simplest solutions of the Einstein equation are the ones most relevant for experiment and observation. The static, spherically symmetric Schwarzschild geometry (3.3) describes the solar system experimental tests, spherically symmetric gravitational collapse, and spherical black holes. The exactly homogeneous and isotropic Friedman-Robertson-Walker (FRW) models provide an excellent approximation to the structure and evolution of our universe from the big bang to the distant future. The linearized solutions of the Einstein equation about a flat space background describe detectable gravitational waves. A ‘physics first’ treatment of general relativity that concentrates on the simplest solutions of the Einstein equation is thus immediately relevant for physically realistic and important situations. No Stopping before Some Physics: Students with different levels of experience, preparation, abilities, and preconceptions will take different lengths of time to acquire the basic concepts of general relativity. Beginning with the applications guarantees that, wherever the course ends, students will have gained some understanding of the basic physical phenomena (Table 1) which make general relativity so important and not simply of a mathematical structure which is the prerequisite for a deductive approach. More Concrete, Less Abstract: Beginning with the applications rather than the abstract structure of the theory is easier for some students because it is more concrete. Beginning with the applications is also a surer way of driving home that general relativity is a part of physics whose predictions can be observationally tested and not a branch of mathematics. Fewer Compromises: The analysis of the motion of test particles and light rays in the simplest geometries can be carried out in essentially the same way as it is done in advanced textbooks. No compromises of method or generality are needed. Flexibility in Emphasis: Beginning with applications allows enough time to construct courses with different emphases; for instance on black holes, gravitational waves, cosmology, or experimental tests. Closer to the Rest of the Undergraduate Physics Curriculum: Calculations of the orbits of test particles and light rays to explore curved spacetimes are exercises in mechanics. The symmetries of the simplest important solutions imply conservation laws. These can be used to reduce the calculation of orbits to one-dimensional motion in effective potentials. Even the content of the Einstein equation can be put in this form for simple situations. This allows the intuition and techniques developed in intermediate mechanics to be brought to bear, both for qualitative understanding and quantitative prediction. Conversely, this kind of example serves to extend and reinforce an understanding of mechanics. Indeed, in the author’s experience a few students are surprised to find that mechanics is actually useful for something. The examples discussed in the next section will help illustrate the close connection between calculating geodesics and undergraduate mechanics. Fewer Prerequisites: The close connection with mechanics described above and in the next section means that the only essential prerequisite to a ‘physics first’ exposition of general relativity is an intermediate course in mechanics. Neither quantum mechanics nor electrodynamics are necessary.Some acquaintance with special relativity is useful, but its brief treatment in many first year courses means that it is usually necessary to develop it de novo at the beginning of an introductory course in general relativity. Intermediate mechanics is thus the single essential physics prerequisite. This means that an introductory ‘physics first’ course in general relativity can be accessible to a wider range of physics majors at an earlier stage than courses designed to introduce students to other frontier areas. Closer to Research Frontiers: Beginning with the applications means that students are closer sooner to the contemporary frontiers of astrophysics and particle physics that they can hear about in the seminars, read about in the newspapers, and see on popular television programs. More Opportunities for Undergraduate Participation in Research: The applications of general relativity provide a broad range of topics for students to pursue independent study or even to make research contributions of contemporary interest. More importantly, it is possible to identify problems from the applications that are conceptually and technically accessible to undergraduate physics majors and can be completed in the limited time frame typically available. Problems that involve solving for the behavior of test particles, light rays, and gyroscopes are examples, as are questions involving linear gravitational waves, black holes, and simple cosmological models. The ‘physics first’ approach to teaching general relativity enables undergraduate participation in research because it treats such applications first. Specialized Faculty Not Needed: Learning a subject while teaching it, or learning it better, is a part of every physics instructor’s experience. The absence of previous instruction means that teaching an introductory general relativity course will often be the first exposure to the subject for many faculty. The process of learning by faculty is made easier by a ‘physics first’ approach for the same reasons it is easier for students. A wider variety of faculty will find this approach both familiar and manageable. Specialists in gravitational physics are not necessary. ### V.2 Counterindications No approach to teaching is without its price and the ‘physics first’ approach to introducing general relativity is no exception. The obvious disadvantage is that it does not follow the logically appealing deductive order, although it can get to the same point in the end. A ‘physics first’ approach to introducing general relativity may not be indicated in at least two circumstances: First, when there is enough time in the curriculum and enough student commitment to pursue the deductive order; second, when students already have significant preparation in differential geometry. Even then however it is a possible alternative. A ‘physics first’ approach is probably not indicated when the mathematics is of central interest, as for students concentrating in mathematics, and for physics students who study Einstein’s theory mainly as an introduction to the mathematics of string theory. ## VI Particle Orbits Outside a Spherical Star or Black Hole This section illustrates how the effective potential method developed in typical undergraduate mechanics courses can be applied to important problems in general relativity. Of the several possible illustrations, just one is considered here — the relativistic effects on test particle orbits outside a spherical star or black hole. Applications of this include the precession of perihelia of planets in the solar system and the location of the innermost stable circular orbit (ISCO) in an accretion disk powering an X-ray source. These examples also serve to illustrate how near the calculations of such effects are to starting principles in a ‘physics first’ approach to general relativity. Nothing more than sketches of the calculations are intended. For more detail and the precise meaning of any quantities involved, the reader should consult any standard text<sup>3</sup><sup>3</sup>3We follow, with one minor simplifying exception, the notation in Har03a .. The Schwarzschild geometry specified in (3.3) describes the curved spacetime outside a static, spherically symmetric body of mass $`M`$. This is the geometry outside the Sun to an excellent approximation, and is the geometry outside a spherical black hole. The motion of test particles is specified by the principle of stationary proper time in (3.4). The argument in the square bracket of that equation can be thought of as a Lagrangian $`L(\dot{t},\dot{r},\dot{\theta },\dot{\varphi },r)`$. Lagrange’s equations are the geodesic equations. Time translation invariance implies a conserved energy per unit rest mass $``$ related to $`L/\dot{t}`$ . Spherical symmetry implies a conserved angular momentum per unit rest mass <sup>4</sup><sup>4</sup>4In Har03a , $``$ and $`\mathrm{}`$ are defined as an energy and angular momentum per unit rest energy rather than per unit rest mass as here. With that definition of $``$ both terms on the right hand side of (6.1a) would be divided by $`c^2`$. That is the minor exception alluded to in the previous footnote. $`\mathrm{}`$ for orbits in the equatorial plane which is proportional to $`L/\dot{\varphi }`$. A third integral<sup>5</sup><sup>5</sup>5‘Integral’ here is used in the sense of classical mechanics, not in the sense of the inverse of differentiation. of the motion $`L=1.`$ arises just from the definition of proper time. This integral can be combined with the other to to find an effective energy integral for the radial motion: $$=\frac{1}{2}\left(\frac{dr}{d\tau }\right)^2+V_{\mathrm{eff}}(r)$$ (6.1a) where $$V_{\mathrm{eff}}(r)=\frac{GM}{r}+\frac{\mathrm{}^2}{2r^2}\frac{GM\mathrm{}^2}{c^2r^3}.$$ (6.1b) Here, $`r`$ is the Schwarzschild radial coordinate of the test particle and $`\tau `$ is the proper time along its world line. The mass of the test particle is absent from these expressions. It cancels out because of the equality of gravitational and inertial mass. Eq. (6.1a) has the same form as the energy integral for a Newtonian central force problem. The first two terms in the effective potential (6.1b) have the same form as a Newtonian gravitational potential with a Newtonian centrifugal barrier. The third term provides a general relativistic correction to the Newtonian effective potential. Figure 2 shows its effects. Circular orbits illustrate the importance of these effects. Newtonian gravity permits only one stable circular orbit for each $`\mathrm{}`$. But in general relativity there are two circular orbits for values of $`\mathrm{}`$ such at that used in Figure 2. There is a stable circular orbit at the minimum of the effective potential such as one approximating the orbit of the Earth in its progress around the Sun. In addition there is an unstable circular orbit at the radius of the maximum of the effective potential. The radii of the stable circular orbits are easily found from (6.1b): $$r_{\mathrm{stab}.\mathrm{circ}.}=\frac{\mathrm{}^2}{2GM}\left\{1+\left[112\left(\frac{GM}{c\mathrm{}}\right)^2\right]^{\frac{1}{2}}\right\}.$$ (6.2) For sufficiently small $`\mathrm{}`$ there are no stable circular orbits. That is because the effective potential (6.1b) is everywhere attractive for low $`\mathrm{}`$. In contrast to Newtonian physics, general relativity therefore implies that there is an innermost (smallest $`r`$) stable circular orbit (ISCO) whose radius is $$r_{\mathrm{ISCO}}=6GM/c^2$$ (6.3) which is $`1.5`$ times the characteristic radius of the black hole $`r_s=2GM/c^2`$. The ISCO is important for the astrophysics of black holes. The spectra of X-ray sources exhibit lines whose observed shapes can in principle be used to infer the properties of the black hole engine Ironlines . The shape of the line is determined by several factors but importantly affected by the gravitational redshift. That is maximum for radiation from parts of the accretion disk closest to the black hole, ie from the ISCO. This defines the red end of the observed line<sup>6</sup><sup>6</sup>6Realistic black holes are generally rotating, but the analysis then is not qualitatively different from that for the non-rotating Schwarzschild black hole.. Another important prediction of general relativity derivable from the effective potential is the shape of bound orbits such as those of the planets. The shape of an orbit in the equatorial plane ($`\theta =\pi /2`$) may be specified by giving the azimuthal angle $`\varphi `$ as a function of $`r`$. The orbits close if the total angle $`\mathrm{\Delta }\varphi `$ swept out in the passage away from the inner turning point and back again is $`2\pi `$. This can be calculated from (6.1) and the angular momentum integral $`\mathrm{}=r^2(d\varphi /d\tau )`$. Writing $`d\varphi /dr=(d\varphi /d\tau )/(dr/d\tau )`$ and using these two relations gives an expression for $`d\varphi /dr`$ as a function of $`r`$, $``$, and $`\mathrm{}`$ which can be integrated. The result is $$\mathrm{\Delta }\varphi =2𝑑r(\mathrm{}/r^2)\left[2\left(V_{\mathrm{eff}}(r)\right)\right]^{1/2}$$ (6.4) where the integral is from the radius of the inner turning point to the outer one. When the relativistic term in (6.1b) is absent, $`\mathrm{\Delta }\varphi =2\pi `$ for all $``$ and $`\mathrm{}`$. That is the closing of the Keplerian ellipses of Newtonian mechanics. The relativistic correction to $`V_{\mathrm{eff}}`$ makes the orbit precess by small amount $`\delta \varphi =\mathrm{\Delta }\varphi 2\pi `$ on each pass. To lowest order in $`1/c^2`$ this is $$\delta \varphi =6\pi \left(\frac{GM}{c\mathrm{}}\right)^2.$$ (6.5) In the solar system the precession is largest for Mercury but still only $`43^{\prime \prime }`$ per century. The confirmation of that prediction Sha90 is an important test of general relativity. The purpose of this section was not to explain or even emphasize the two effects of general relativity on particle orbits that were described here. Rather it was to show three things. First, that a standard technique developed in undergraduate mechanics can be used to calculate important predictions of general relativity. The second purpose was to show how close these important applications can come to starting principles in a ‘physics first’ approach to teaching general relativity. Introduce the Schwarzschild geometry (3.3), use the principle of stationary proper time (3.4) to find the geodesic equations or their integrals (6.1), use the effective potential method to qualitatively and quantitatively understand important properties of the orbits e.g. (6.3) and (6.4). That is just three steps from the basic ideas of metric and geodesics to important applications. Third, both of the applications treated here can be immediately related to contemporary observation and experiment. Particle orbits in the Schwarzschild geometry are not the only important problems in which the effective potential method is useful. Motion in the geometry of a rotating black hole, the motion of test light rays in the Schwarzschild geometry, and the evolution of the Friedman-Robertson-Walker cosmological models are further examples where it can be usefully applied. ## VII Conclusion A one quarter ($``$ 28 lectures) ‘physics first’ course in general relativity has been a standard junior/senior elective for physics majors at the University of California, Santa Barbara for approximately thirty years. In the limited span of a quarter the author is usually able to review special relativity, motivate gravity as geometry, derive the orbits in the Schwarzschild geometry in detail, describe the experimental tests, introduce black holes, and develop the Friedman-Robertson-Walker cosmological models. A semester provides more opportunities. At Santa Barbara this course is routinely taught by faculty from many different areas of physics — general relativity of course, but also elementary particle physics and astrophysics. It has been taught by both theorists and experimentalists. For some of these colleagues teaching this course was their first serious experience with general relativity. They usually report that they were successful and enjoyed it. In the author’s experience students are excited by general relativity and motivated to pursue it. Often it is their first experience with a subject directly relevant to current research. It is one of the few contemporary subjects that can be taught without quantum mechanics or electromagnetism. The author has written a text based on the ‘physics first’ approach Har03a which comes with a solutions manual available to instructors for the approximately 400 problems of graded levels of difficulty. ‘Physics first’ is not the only way of introducing general relativity to undergraduate physics majors, but it works. ###### Acknowledgements. Thanks are due to Bill Hiscock, Ted Jacobson, Richard Price, and Francesc Roig for critical readings of the manuscript. Preparation of this paper was supported in part by NSF grant PHY02-44764.
warning/0506/gr-qc0506059.html	ar5iv	text	# LISA Data Analysis using MCMC methods ## I Introduction The LISA observatory lppa has incredible science potential, but that potential can only be fully realized by employing advanced data analysis techniques. LISA will explore the low frequency portion of the gravitational wave spectrum, which is thought to be home to a vast number of sources. Since gravitational wave sources typically evolve on timescales that are long compared to the gravitational wave period, individual low frequency sources will be “on” for large fractions of the nominal three year LISA mission lifetime. Moreover, unlike a traditional telescope, LISA can not be pointed at a particular point on the sky. The upshot is that the LISA data stream will contain the signals from tens of thousands of individual sources, and ways must be found to isolate individual voices from the crowd. This “Cocktail Party Problem” is the central issue in LISA data analysis. The types of sources LISA is expected to detect include galactic and extra-galactic compact stellar binaries, super massive black hole binaries, and extreme mass ratio inspirals of compact stars into supermassive black holes (EMRIs). Other potential sources include intermediate mass black hole binaries, cosmic strings, and a cosmic gravitational wave background produced by processes in the early universe. In the case of compact stellar binaries evans ; lip ; hils ; hils2 ; gils and EMRIs cutbar ; emri , the number of sources is likely to be so large that it will be impossible to resolve all the sources individually, so that there will be a residual signal that is variously referred to as a confusion limited background or confusion noise. It is important that this confusion noise be made as small as possible so as not to hinder the detection of other high value targets. Several estimates of the confusion noise level have been made hils ; hils2 ; gils ; sterl ; seth ; bc , and they all suggest that unresolved signals will be the dominant source of low frequency noise for LISA. However, these estimates are based on assumptions about the efficacy of the data analysis algorithms that will be used to identify and regress sources from the LISA data stream, and it is unclear at present how reasonable these assumptions might be. Indeed, the very notion that one can first clean the data stream of one type of signal before moving on to search for other targets is suspect as the gravitational wave signals from different sources are not orthogonal. For example, when the signal from a supermassive black hole binary sweeps past the signal from a white dwarf binary of period $`T`$, the two signals will have significant overlap for a time interval equal to the geometric mean of $`T`$ and $`t_c`$, where $`t_c`$ is the time remaining before the black holes merge. Thus, by a process dubbed “the white dwarf transform,” it is possible to decompose the signal from a supermassive black hole binary into signals from a collection of white dwarf binaries. As described in §II, optimal filtering of the LISA data would require the construction of a filter bank that described the signals from every source that contributes to the data stream. In principle one could construct a vast template bank describing all possible sources and look for the best match with the data. In practice the enormous size of the search space and the presence of unmodeled sources renders this direct approach impractical. Possible alternatives to a full template based search include iterative refinement of a source-by-source search, ergodic exploration of the parameter space using Markov Chain Monte Carlo (MCMC) algorithms , Darwinian optimization by genetic algorithms, and global iterative refinement using the Maximum Entropy Method (MEM). Each approach has its strengths and weakness, and at this stage it is not obvious which approach will prove superior. Here we apply the popular Markov Chain Monte Carlo metro ; haste method to simulated LISA data. This is not the first time that MCMC methods have been applied to gravitational wave data analysis, but it is first outing with realistic simulated LISA data. Our simulated data streams contain the signals from multiple galactic binaries. Previously, MCMC methods have been used to study the extraction of coalescing binary christ and spinning neutron star woan signals from terrestrial interferometers. More recently, MCMC methods have been applied to a simplified toy problem woan2 that shares some of the features of the LISA cocktail party problem. These studies have shown that MCMC methods hold considerable promise for gravitational wave data analysis, and offer many advantages over the standard template grid searches. For example, the EMRI data analysis problem cutbar ; emri is often cited as the greatest challenge facing LISA science. Neglecting the spin of the smaller body yields a 14 dimensional parameter space, which would require $`10^{40}`$ templates to explore in a grid based search emri . This huge computational cost arises because grid based searches scale geometrically with the parameter space dimension $`D`$. In contrast, the computational cost of MCMC based searches scale linearly with the $`D`$. In fields such as finance, MCMC methods are routinely applied to problems with $`D>1000`$, making the LISA EMRI problem seem trivial in comparison. A Google search on “Markov Chain Monte Carlo” returns almost 250,000 results, and a quick scan of these pages demonstrates the wide range of fields where MCMC methods are routinely used. We found it amusing that one of the Google search results is a link to the PageRank page MCMC algorithm that powers the Google search engine. The structure of the paper follows the development sequence we took to arrive at a fast and robust MCMC algorithm. In §II we outline the LISA data analysis problem and the particular challenges posed by the galactic background. A basic MCMC algorithm is introduced in §III and applied to a full 7 parameter search for a single galactic binary. A generalized multi-channel, multi-source F-statistic for reducing the search space from $`D=7N`$ to $`D=3N`$ is described in §IV. The performance of a basic MCMC algorithm that uses the F-statistic is studied in §V and a number of problems with this simple approach are identified. A more advanced mixed MCMC algorithm that incorporates simulated annealing is introduced in §VI and is successfully applied to multi-source searches. The issue of model selection is addressed in §VII, and approximate Bayes factor are calculated by super-cooling the Markov Chains to extract maximum likelihood estimates. We conclude with a discussion of future refinements and extensions of our approach in §VIII. ## II The Cocktail Party Problem Space based detectors such as LISA are able to return several interferometer outputs aet . The strains registered in the interferometer in response to a gravitational wave pick up modulations due to the motion of the detector. The orbital motion introduces amplitude, frequency, and phase modulation into the observed gravitational wave signal. The amplitude modulation results from the detector’s antenna pattern being swept across the sky, the frequency modulation is due to the Doppler shift from the relative motion of the detector and source, and the phase modulation results from the detector’s varying response to the two gravitational wave polarizations cc ; cr . These modulations encode information about the location of the source. The modulations spread a monochromatic signal over a bandwidth $`\mathrm{\Delta }f(9+6(f/\mathrm{mHz})\mathrm{sin}\theta )f_m`$, where $`\theta `$ is the co-latitude of the source and $`f_m=1/\mathrm{year}`$ is the modulation frequency. In the low frequency limit, where the wavelengths are large compared to the armlengths of the detector, the interferometer outputs $`s_\alpha (t)`$ can be combined to simulate the response of two independent 90 degree interferometers, $`s_I(t)`$ and $`s_{II}(t)`$, rotated by 45 degrees with respect to each other cc ; tom . This allows LISA to measure both polarizations of the gravitational wave simultaneously. A third combination of signals in the low frequency limit yields the symmetric Sagnac variable aet , which is insensitive to gravitational waves and can be used to monitor the instrument noise. When the wavelengths of the gravitational waves become comparable to the size of the detector, which for LISA corresponds to frequencies above 10 mHz, the interferometry signals can be combined to give three independent time series with comparable sensitivities tom . The output of each LISA data stream can be written as $$s_\alpha (t)=h_\alpha (t,\stackrel{}{\lambda })+n_\alpha (t)=\underset{i=1}{\overset{N}{}}h_\alpha ^i(t,\stackrel{}{\lambda }_i)+n_\alpha (t).$$ (1) Here $`h_\alpha ^i(t,\stackrel{}{\lambda }_i)`$ describes the response registered in detector channel $`\alpha `$ to a source with parameters $`\stackrel{}{\lambda }_i`$. The quantity $`h_\alpha (t,\stackrel{}{\lambda })`$ denotes the combined response to a collection of $`N`$ sources with total parameter vector $`\stackrel{}{\lambda }=_i\stackrel{}{\lambda }_i`$ and $`n_\alpha (t)`$ denotes the instrument noise in channel $`\alpha `$. Extracting the parameters of each individual source from the combined response to all sources defines the LISA cocktail party problem. In practice it will be impossible to resolve all of the millions of signals that contribute to the LISA data streams. For one, there will not be enough bits of information in the entire LISA data archive to describe all $`N`$ sources in the Universe with signals that fall within the LISA band. Moreover, most sources will produce signals that are well below the instrument noise level, and even after optimal filtering most of these sources will have signal to noise ratios below one. A more reasonable goal might be to provide estimates for the parameters describing each of the $`N^{}`$ sources that have integrated signal to noise ratios (SNR) above some threshold (such as $`\mathrm{SNR}>5`$), where it is now understood that the noise includes the instrument noise, residuals from the regression of bright sources, and the signals from unresolved sources. While the noise will be neither stationary nor Gaussian, it is not unreasonable to hope that the departures from Gaussianity and stationarity will be mild. It is well know that matched filtering is the optimal linear signal processing technique for signals with stationary Gaussian noise helstrom ; wz . Matched filtering is used extensively in all fields of science, and is a popular data analysis technique in ground based gravitational wave astronomy kip ; bernie ; sathya1 ; sathya2 ; curt1 ; bala1 ; sathya3 ; ap1 ; eric1 ; sathya4 ; ben1 ; ben2 . Switching to the Fourier domain, the signal can be written as $`\stackrel{~}{s}_\alpha (f)=\stackrel{~}{h}_\alpha (f,\stackrel{}{\lambda }^{})+\stackrel{~}{n}_\alpha (f)`$, where $`\stackrel{~}{n}_\alpha (f)`$ includes instrument noise and confusion noise, and the signals are described by parameters $`\stackrel{}{\lambda }^{}`$. Using the standard noise weighted inner product for the independent data channels over a finite observation time $`T`$, $$(a\|b)=\frac{2}{T}\underset{\alpha }{}\underset{f}{}\frac{\stackrel{~}{a}_\alpha ^{}(f)\stackrel{~}{b}_\alpha (f)+\stackrel{~}{a}_\alpha (f)\stackrel{~}{b}_\alpha ^{}(f)}{S_n^\alpha (f)},$$ (2) a Wiener filter statistic can be defined: $$\rho (\stackrel{}{\lambda })=\frac{(s\|h(\stackrel{}{\lambda }))}{\sqrt{(h(\stackrel{}{\lambda })\|h(\stackrel{}{\lambda }))}}.$$ (3) The noise spectral density $`S_n(f)`$ is given in terms of the autocorrelation of the noise $$n(f)n^{}(f^{})=\frac{T}{2}\delta _{ff^{}}S_n(f).$$ (4) Here and elsewhere angle brackets $``$ denote an expectation value. An estimate for the source parameters $`\stackrel{}{\lambda }^{}`$ can be found by maximizing $`\rho (\stackrel{}{\lambda })`$. If the noise is Gaussian and a signal is present, $`\rho (\stackrel{}{\lambda })`$ will be Gaussian distributed with unit variance and mean equal to the integrated signal to noise ratio $$\mathrm{SNR}=\rho (\stackrel{}{\lambda }^{})=\sqrt{(h(\stackrel{}{\lambda }^{})\|h(\stackrel{}{\lambda }^{}))}.$$ (5) The optimal filter for the LISA signal (1) is a matched template describing all $`N^{}`$ resolvable sources. The number of parameters $`d_i`$ required to describe a source ranges from 7 for a slowly evolving circular galactic binary to 17 for a massive black hole binary. A reasonable estimate seth for $`N^{}`$ is around $`10^4`$, so the full parameter space has dimension $`D=_id_i10^5`$. Since the number of templates required to uniformly cover a parameter space grows exponentially with $`D`$, a grid based search using the full optimal filter is out of the question. Clearly an alternative approach has to be found. Moreover, the number of resolvable sources $`N^{}`$ is not known a priori, so some stopping criteria must be found to avoid over-fitting the data. Existing approaches to the LISA cocktail party problem employ iterative schemes. The first such approach was dubbed “gCLEAN” gclean due to its similarity with the “CLEAN” Hoegbom algorithm that is used for astronomical image reconstruction. The “gCLEAN” procedure identifies and records the brightest source that remains in the data stream, then subtracts a small amount of this source. The procedure is iterated until a prescribed residual is reached, at which time the individual sources are reconstructed from the subtraction record. A much faster iterative approach dubbed “Slice & Dice” slicedice was recently proposed that proceeds by identifying and fully subtracting the brightest source that remains in the data stream. A global least squares re-fit to all the current list of sources is then performed, and the new parameter record is used to produce a regressed data stream for the next iteration. Bayes factors are used to provide a stopping criteria. There is always the danger with iterative approaches that the procedure “gets off on the wrong foot,” and is unable to find its way back to the optimal solution. This can happen when two signals have a high degree of overlap. A very different approach to the LISA source confusion problem is to solve for all sources simultaneously using ergodic sampling techniques. Markov Chain Monte Carlo (MCMC) gilks ; gamer is a method for estimating the posterior distribution, $`p(\stackrel{}{\lambda }\|s)`$, that can be used with very large parameter spaces. The method is now in widespread use in many fields, and is starting to be used by astronomers and cosmologists. One of the advantages of MCMC is that it combines detection, parameter estimation, and the calculation of confidence intervals in one procedure, as everything one can ask about a model is contained in $`p(\stackrel{}{\lambda }\|s)`$. Another nice feature of MCMC is that there are implementations that allow the number of parameters in the model to be variable, with built in penalties for using too many parameters in the fit. In an MCMC approach, parameter estimates from Wiener matched filtering are replaced by the Bayes estimator davis $$\lambda _\mathrm{B}^i(s)=\lambda ^ip(\stackrel{}{\lambda }\|s)𝑑\stackrel{}{\lambda },$$ (6) which requires knowledge of $`p(\stackrel{}{\lambda }\|s)`$ \- the posterior distribution of $`\stackrel{}{\lambda }`$ (i.e. the distribution of $`\stackrel{}{\lambda }`$ conditioned on the data $`s`$). By Bayes theorem, the posterior distribution is related to the prior distribution $`p(\stackrel{}{\lambda })`$ and the likelihood $`p(s\|\stackrel{}{\lambda })`$ by $$p(\stackrel{}{\lambda }\|s)=\frac{p(\stackrel{}{\lambda })p(s\|\stackrel{}{\lambda })}{p(\stackrel{}{\lambda ^{}})p(s\|\stackrel{}{\lambda ^{}})𝑑\stackrel{}{\lambda ^{}}}.$$ (7) Until recently the Bayes estimator was little used in practical applications as the integrals appearing in (6) and (7) are often analytically intractable. The traditional solution has been to use approximations to the Bayes estimator, such as the maximum likelihood estimator described below, however advances in the Markov Chain Monte Carlo technique allow direct numerical estimates to be made. When the noise $`n(t)`$ is a normal process with zero mean, the likelihood is given by sam $$p(s\|\stackrel{}{\lambda })=C\mathrm{exp}\left[\frac{1}{2}\left((sh(\stackrel{}{\lambda }))\|(sh(\stackrel{}{\lambda }))\right)\right],$$ (8) where the normalization constant $`C`$ is independent of $`s`$. In the large SNR limit the Bayes estimator can be approximated by finding the dominant mode of the posterior distribution, $`p(\stackrel{}{\lambda }\|s)`$, which Finn sam and Cutler & Flannagancurt1 refer to as a maximum likelihood estimator. Other authors kro ; fstat define the maximum likelihood estimator to be the value of $`\stackrel{}{\lambda }`$ that maximizes the likelihood, $`p(s\|\stackrel{}{\lambda })`$. The former has the advantage of incorporating prior information, but the disadvantage of not being invariant under parameter space coordinate transformations. The latter definition corresponds to the standard definition used by most statisticians, and while it does not take into account prior information, it is coordinate invariant. The two definitions give the same result for uniform priors, and very similar results in most cases (the exception being where the priors have a large gradient at maximum likelihood). The standard definition of the likelihood yields an estimator that is identical to Wiener matched filteringeche . Absorbing normalization factors by adopting the inverted relative likelihood $`(\stackrel{}{\lambda })=p(s\|0)/p(s\|\stackrel{}{\lambda })`$, we have $$\mathrm{log}(\stackrel{}{\lambda })=(s\|h(\stackrel{}{\lambda }))\frac{1}{2}(h(\stackrel{}{\lambda })\|h(\stackrel{}{\lambda })).$$ (9) In the gravitational wave literature the quantity $`\mathrm{log}(\stackrel{}{\lambda })`$ is usually referred to as the log likelihood, despite the inversion and rescaling. Note that $$\mathrm{log}(\stackrel{}{\lambda }^{})=\frac{1}{2}\rho (\stackrel{}{\lambda }^{})^2=\frac{1}{2}\mathrm{SNR}^2.$$ (10) The maximum likelihood estimator (MLE), $`\stackrel{}{\lambda }_{\mathrm{ML}}`$, is found by solving the coupled set of equations $`\mathrm{log}/\lambda ^i=0`$. Parameter uncertainties can be estimated from the negative Hessian of $`\mathrm{log}`$, which yields the Fisher Information Matrix $$\mathrm{\Gamma }_{ij}(\stackrel{}{\lambda })=\frac{^2\mathrm{log}(\stackrel{}{\lambda })}{\lambda ^i\lambda ^j}=(h_{,i}\|h_{,j}).$$ (11) In the large SNR limit the MLE can be found by writing $`\stackrel{}{\lambda }=\stackrel{}{\lambda }^{}+\mathrm{\Delta }\stackrel{}{\lambda }`$ and Taylor expanding (9). Setting $`\mathrm{log}/\mathrm{\Delta }\lambda ^i=0`$ yields the lowest order solution $$\lambda _{\mathrm{ML}}^i=\lambda _{}^{}{}_{}{}^{i}+\mathrm{\Delta }\lambda ^i=\lambda _{}^{}{}_{}{}^{i}+\mathrm{\Gamma }^{ij}(\stackrel{}{\lambda ^{}})(n\|h_{,j}).$$ (12) The expectation value of the maximum of the log likelihood is then $$\mathrm{log}(\stackrel{}{\lambda }_{\mathrm{ML}})=\frac{\mathrm{SNR}^2+D}{2}.$$ (13) This value exceeds that found in (10) by an amount that depends on the total number of parameters used in the fit, $`D`$, reflecting the fact that models with more parameters generally give better fits to the data. Deciding how many parameters to allow in the fit is an important issue in LISA data analysis as the number of resolvable sources is not known a priori. This issue does not usually arise for ground based gravitational wave detectors as most high frequency gravitational wave sources are transient. The relevant question there is whether or not a gravitational wave signal is present in a section of the data stream, and this question can be dealt with by the Neyman-Pearson test or other similar tests that use thresholds on the likelihood $``$ that are related to the false alarm and false dismissal rates. Demanding that $`>1`$ \- so it is more likely that a signal is present than not - and setting a detection threshold of $`\rho =5`$ yields a false alarm probability of 0.006 and a detection probability of 0.994 (if the noise is stationary and Gaussian). A simple acceptance threshold of $`\rho =5`$ for each individual signal used to fit the LISA data would help restrict the total number of parameters in the fit, however there are better criteria that can be employed. The simplest is related to the Neyman-Pearson test and compares the likelihoods of models with different numbers of parameters. For nested models this ratio has an approximately chi squared distribution which allows the significance of adding extra parameters to be determined from standard statistical tables. A better approach is to compute the Bayes factor, $$B_{XY}=\frac{p_X(s)}{p_Y(s)},$$ (14) which gives the relative weight of evidence for models $`X`$ and $`Y`$ in terms of the ratio of marginal likelihoods $$p_X(s)=p(s\|\stackrel{}{\lambda },X)p(\stackrel{}{\lambda },X)𝑑\stackrel{}{\lambda }.$$ (15) Here $`p(s\|\stackrel{}{\lambda },X)`$ is the likelihood distribution for model $`X`$ and $`p(\stackrel{}{\lambda },X)`$ is the prior distribution for model $`X`$. The difficulty with this approach is that the integral in (15) is hard to calculate, though estimates can be made using the Laplace approximation or the Bayesian Information Criterion (BIC) schwarz . The Laplace approximation is based on the method of steepest descents, and for uniform priors yields $$p_X(s)p(s\|\stackrel{}{\lambda }_{\mathrm{ML}},X)\left(\frac{\mathrm{\Delta }V_X}{V_X}\right),$$ (16) where $`p(s\|\stackrel{}{\lambda }_{\mathrm{ML}},X)`$ is the maximum likelihood for the model, $`V_X`$ is the volume of the model’s parameter space, and $`\mathrm{\Delta }V_X`$ is the volume of the uncertainty ellipsoid (estimated using the Fisher matrix). Models with more parameters generally provide a better fit to the data and a higher maximum likelihood, but they get penalized by the $`\mathrm{\Delta }V_X/V_X`$ term which acts as a built in Occam’s razor. ## III Markov Chain Monte Carlo We begin by implementing a basic MCMC search for galactic binaries that searches over the full $`D=7N`$ dimensional parameter space using the Metropolis-Hastings haste algorithm. The idea is to generate a set of samples, $`\{\stackrel{}{x}\}`$, that correspond to draws from the posterior distribution, $`p(\stackrel{}{\lambda }\|s)`$. To do this we start at a randomly chosen point $`\stackrel{}{x}`$ and generate a Markov chain according to the following algorithm: Using a proposal distribution $`q(\|\stackrel{}{x})`$, draw a new point $`\stackrel{}{y}`$. Evaluate the Hastings ratio $$H=\frac{p(\stackrel{}{y})p(s\|\stackrel{}{y})q(\stackrel{}{x}\|\stackrel{}{y})}{p(\stackrel{}{x})p(s\|\stackrel{}{x})q(\stackrel{}{y}\|\stackrel{}{x})}.$$ (17) Accept the candidate point $`\stackrel{}{y}`$ with probability $`\alpha =\mathrm{min}(1,H)`$, otherwise remain at the current state $`\stackrel{}{x}`$ (Metropolis rejection metro ). Remarkably, this sampling scheme produces a Markov chain with a stationary distribution equal to the posterior distribution of interest, $`p(\stackrel{}{\lambda }\|s)`$, regardless of the choice of proposal distribution gilks . A concise introduction to MCMC methods can be found in the review paper by Andrieu et al mcmc\_hist . On the other hand, a poor choice of the proposal distribution will result in the algorithm taking a very long time to converge to the stationary distribution (known as the burn-in time). Elements of the Markov chain produced during the burn-in phase have to be discarded as they do not represent the stationary distribution. When dealing with large parameter spaces the burn-in time can be very long if poor techniques are used. For example, the Metropolis sampler, which uses symmetric proposal distributions, explores the parameter space with an efficiency of at most $`0.3/D`$, making it a poor choice for high dimension searches. Regardless of the sampling scheme, the mixing of the Markov chain can be inhibited by the presence of strongly correlated parameters. Correlated parameters can be dealt with by making a local coordinate transformation at $`\stackrel{}{x}`$ to a new set of coordinates that diagonalises the Fisher matrix, $`\mathrm{\Gamma }_{ij}(\stackrel{}{x})`$. We tried a number of proposal distributions and update schemes to search for a single galactic binary. The results were very disappointing. Bold proposals that attempted large jumps had a very poor acceptance rate, while timid proposals that attempted small jumps had a good acceptance rate, but they explored the parameter space very slowly, and got stuck at local modes of the posterior. Lorentzian proposal distributions fared the best as their heavy tails and concentrated peaks lead to a mixture of bold and timid jumps, but the burn in times were still very long and the subsequent mixing of the chain was torpid. The MCMC literature is full of similar examples of slow exploration of large parameter spaces, and a host of schemes have been suggested to speed up the burn-in. Many of the accelerated algorithms use adaptation to tune the proposal distribution. This violates the Markov nature of the chain as the updates depend on the history of the chain. More complicated adaptive algorithms have been invented that restore the Markov property by using additional Metropolis rejection steps. The popular Delayed Rejection Method dr and Reversible Jump Method rj are examples of adaptive MCMC algorithms. A simpler approach is to use a non-Markov scheme during burn-in, such as adaptation or simulated annealing, then transition to a Markov scheme after burn-in. Since the burn-in portion of the chain is discarded, it does not matter if the MCMC rules are broken (the burn-in phase is more like Las Vegas than Monte Carlo). Before resorting to complex acceleration schemes we tried a much simpler approach that proved to be very successful. When using the Metropolis-Hastings algorithm there is no reason to restrict the updates to a single proposal distribution. For example, every update could use a different proposal distribution so long as the choice of distribution is not based on the history of the chain. The proposal distributions to be used at each update can be chosen at random, or they can be applied in a fixed sequence. Our experience with single proposal distributions suggested that a scheme that combined a very bold proposal with a very timid proposal would lead to fast burn-in and efficient mixing. For the bold proposal we chose a uniform distribution for each of the source parameters $`\stackrel{}{\lambda }(A,f,\theta ,\varphi ,\psi ,\iota ,\phi _0)`$. Here $`A`$ is the amplitude, $`f`$ is the gravitational wave frequency, $`\theta `$ and $`\varphi `$ are the ecliptic co-latitude and longitude, $`\psi `$ is the polarization angle, $`\iota `$ is the inclination of the orbital plane, and $`\phi _0`$ is the orbital phase at some fiducial time. The amplitudes were restricted to the range $`A[10^{23},10^{21}]`$ and the frequencies were restricted to lie within the range of the data snippet $`f[0.999995,1.003164]`$ mHz (the data snippet contained 100 frequency bins of width $`\mathrm{\Delta }f=1/\mathrm{year}`$). A better choice would have been to use a cosine distribution for the co-latitude $`\theta `$ and inclination $`\iota `$, but the choice is not particularly important. When multiple sources were present each source was updated separately during the bold proposal stage. For the timid proposal we used a normal distribution for each eigendirection of the Fisher matrix, $`\mathrm{\Gamma }_{ij}(\stackrel{}{x})`$. The standard deviation $`\sigma _{\widehat{k}}`$ for each eigendirection $`k`$ was set equal to $`\sigma _{\widehat{k}}=1/\sqrt{\alpha _{\widehat{k}}D}`$, where $`\alpha _{\widehat{k}}`$ is the corresponding eigenvalue of $`\mathrm{\Gamma }_{ij}(\stackrel{}{x})`$, and $`D=7N`$ is the search dimension. The factor of $`1/\sqrt{D}`$ ensures a healthy acceptance rate as the typical total jump is then $`1\sigma `$. All $`N`$ sources were updated simultaneously during the timid proposal stage. Note that the timid proposal distributions are not symmetric since $`\mathrm{\Gamma }_{ij}(\stackrel{}{x})\mathrm{\Gamma }_{ij}(\stackrel{}{y})`$. One set of bold proposals (one for each source) was followed by ten timid proposals in a repeating cycle. The ratio of the number of bold to timid proposals impacted the burn-in times and the final mixing rate, but ratios anywhere from 1:1 to 1:100 worked well. We used uniform priors, $`p(\stackrel{}{x})=\mathrm{const}.`$, for all the parameters, though once again a cosine distribution would have been better for $`\theta `$ and $`\iota `$. Two independent LISA data channels were simulated directly in the frequency domain using the method described in Ref. seth , with the sources chosen at random using the same uniform distributions employed by the bold proposal. The data covers 1 year of observations, and the data snippet contains 100 frequency bins (of width $`1/\mathrm{year}`$). The instrument noise was assumed to be stationary and Gaussian, with position noise spectral density $`S_n^{\mathrm{pos}}=4\times 10^{22}\mathrm{m}^2\mathrm{Hz}^1`$ and acceleration noise spectral density $`S_n^{\mathrm{accel}}=9\times 10^{30}\mathrm{m}^2\mathrm{s}^4\mathrm{Hz}^1`$. Table 1 summarizes the results of one MCMC run using a model with one source to search for a single source in the data snippet. Burn-in lasted $`2000`$ iterations, and post burn-in the chain was run for $`10^6`$ iterations with a proposal acceptance rate of $`77\%`$ (the full run took 20 minutes on a Mac G5 2 GHz processor). The chain was used to calculate means and variances for all the parameters. The parameter uncertainty estimates extracted from the MCMC output are compared to the Fisher matrix estimates evaluated at the mean values of the parameters. The source had true $`\mathrm{SNR}=12.9`$, and MCMC recovered $`\mathrm{SNR}=10.7`$. Histograms of the posterior parameter distributions are shown in Figure 1, where they are compared to the Gaussian approximation to the posterior given by the Fisher matrix. The agreement is impressive, especially considering that the bandwidth of the source is roughly 10 frequency bins, so there are very few noise samples to work with. Similar results were found for other MCMC runs on the same source, and for MCMC runs with other sources. Typical burn-in times were of order 3000 iterations, and the proposal acceptance rate was around $`75\%`$. The algorithm was run successfully on two and three source searches (the model dimension was chosen to match the number of sources in each instance), but on occasions the chain would get stuck at a local mode of the posterior for a large number of iterations. Before attempting to cure this problem with a more refined MCMC algorithm, we decided to eliminate the extrinsic parameters $`A,\iota ,\psi ,\epsilon _0`$ from the search by using a multi-filter generalized F-statistic. This reduces the search dimension to $`D=3N`$, with the added benefit that the projection onto the $`(f,\theta ,\varphi )`$ sub-space yields a softer target for the MCMC search. ## IV Generalized F-Statistic The F-statistic was originally introduced fstat in the context of ground based searches for gravitational wave signals from rotating Neutron stars. The F-statistic has since been used to search for monochromatic galactic binaries using simulated LISA data flisa ; slicedice . By using multiple linear filters, the F-statistic is able to automatically extremize the log likelihood over extrinsic parameters, thus reducing the dimension of the search space (the parameter space dimension remains the same). In the low-frequency limit the LISA response to a gravitational wave with polarization content $`h_+(t)`$, $`h_\times (t)`$ can be written as $$h(t)=h_+(t)F^+(t)+h_\times (t)F^\times (t),$$ (18) where $`F^+(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{cos}2\psi D^+(t)\mathrm{sin}2\psi D^\times (t)\right)`$ $`F^\times (t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{sin}2\psi D^+(t)+\mathrm{cos}2\psi D^\times (t)\right)`$ (19) The detector pattern functions $`D^+(t)`$ and $`D^\times (t)`$ are given in equations (36) and (37) of Ref.rigad . To leading post-Newtonian order a slowly evolving, circular binary has polarization components $`h_+(t)`$ $`=`$ $`A(1+\mathrm{cos}^2\iota )\mathrm{cos}(\mathrm{\Phi }(t)+\phi _0)`$ $`h_\times (t)`$ $`=`$ $`2A\mathrm{cos}\iota \mathrm{sin}(\mathrm{\Phi }(t)+\phi _0).`$ (20) The gravitational wave phase $$\mathrm{\Phi }(t;f,\theta ,\varphi )=2\pi ft+2\pi f\mathrm{AU}\mathrm{sin}\theta \mathrm{cos}(2\pi f_mt\varphi ),$$ (21) couples the sky location and the frequency through the term that depends on the radius of LISA’s orbit, 1 AU, and the orbital modulation frequency, $`f_m=1/\mathrm{year}`$. The gravitational wave amplitude, $`A`$, is effectively constant for the low frequency galactic sources we are considering. Using these expressions (18) can be written as $$h(t)=\underset{i=1}{\overset{4}{}}a_i(A,\psi ,\iota ,\phi _0)A^i(t;f,\theta ,\varphi ),$$ (22) where the time-independent amplitudes $`a_i`$ are given by $`a_1`$ $`=`$ $`{\displaystyle \frac{A}{2}}\left((1+\mathrm{cos}^2\iota )\mathrm{cos}\phi _0\mathrm{cos}2\psi 2\mathrm{cos}\iota \mathrm{sin}\phi _0\mathrm{sin}2\psi \right),`$ $`a_2`$ $`=`$ $`{\displaystyle \frac{A}{2}}\left(2\mathrm{cos}\iota \mathrm{sin}\phi _0\mathrm{cos}2\psi +(1+\mathrm{cos}^2\iota )\mathrm{cos}\phi _0\mathrm{sin}2\psi \right),`$ $`a_3`$ $`=`$ $`{\displaystyle \frac{A}{2}}\left(2\mathrm{cos}\iota \mathrm{cos}\phi _0\mathrm{sin}2\psi +(1+\mathrm{cos}^2\iota )\mathrm{sin}\phi _0\mathrm{cos}2\psi \right),`$ $`a_4`$ $`=`$ $`{\displaystyle \frac{A}{2}}\left((1+\mathrm{cos}^2\iota )\mathrm{sin}\phi _0\mathrm{sin}2\psi 2\mathrm{cos}\iota \mathrm{cos}\phi _0\mathrm{cos}2\psi \right),`$ and the time-dependent functions $`A^i(t)`$ are given by $`A^1(t)`$ $`=`$ $`D^+(t;\theta ,\varphi )\mathrm{cos}\mathrm{\Phi }(t;f,\theta ,\varphi )`$ $`A^2(t)`$ $`=`$ $`D^\times (t;\theta ,\varphi )\mathrm{cos}\mathrm{\Phi }(t;f,\theta ,\varphi )`$ $`A^3(t)`$ $`=`$ $`D^+(t;\theta ,\varphi )\mathrm{sin}\mathrm{\Phi }(t;f,\theta ,\varphi )`$ $`A^4(t)`$ $`=`$ $`D^\times (t;\theta ,\varphi )\mathrm{sin}\mathrm{\Phi }(t;f,\theta ,\varphi ).`$ (24) Defining the four constants $`N^i=(s\|A^i)`$ and using (22) yields a solution for the amplitudes $`a_i`$: $$a_i=(M^1)_{ij}N^j,$$ (25) where $`M^{ij}=(A^i\|A^j)`$. The output of the four linear filters, $`N^i`$, and the $`4\times 4`$ matrix $`M^{ij}`$ can be calculated using the same fast Fourier space techniques seth used to generate the full waveforms. Substituting (22) and (25) into expression (9) for the log likelihood yields the F-statistic $$=\mathrm{log}=\frac{1}{2}(M^1)_{ij}N^iN^j.$$ (26) The F-statistic automatically maximizes the log likelihood over the extrinsic parameters $`A,\iota ,\psi `$ and $`\phi _0`$, and reduces the search to the sub-space spanned by $`f,\theta `$ and $`\varphi `$. The extrinsic parameters can be recovered from the $`a_i`$’s via $`A`$ $`=`$ $`{\displaystyle \frac{A_++\sqrt{A_+^2A_\times ^2}}{2}}`$ $`\psi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{arctan}\left({\displaystyle \frac{A_+a_4A_\times a_1}{(A_\times a_2+A_+a_3)}}\right)`$ $`\iota `$ $`=`$ $`\mathrm{arccos}\left({\displaystyle \frac{A_\times }{A_++\sqrt{A_+^2A_\times ^2}}}\right)`$ $`\phi _0`$ $`=`$ $`\mathrm{arctan}\left({\displaystyle \frac{c(A_+a_4A_\times a_1)}{c(A_\times a_2+A_+a_3)}}\right)`$ (27) where $`A_+`$ $`=`$ $`\sqrt{(a_1+a_4)^2+(a_2a_3)^2}`$ $`+\sqrt{(a_1a_4)^2+(a_2+a_3)^2}`$ $`A_\times `$ $`=`$ $`\sqrt{(a_1+a_4)^2+(a_2a_3)^2}`$ $`\sqrt{(a_1a_4)^2+(a_2+a_3)^2}`$ $`c`$ $`=`$ $`\mathrm{sign}(\mathrm{sin}(2\psi )).`$ (28) The preceding description of the F-statistic automatically incorporates the two independent LISA channels through the use of the dual-channel noise weighted inner product $`(a\|b)`$. The basic F-statistic can easily be generalized to handle $`N`$ sources. Writing $`i=4K+l`$, where $`K`$ labels the source and $`l=14`$ labels the four filters for each source, the F-statistic (26) keeps the same form as before, but now there are $`4N`$ linear filters $`N^i`$, and $`M^{ij}`$ is a $`4N\times 4N`$ dimensional matrix. For slowly evolving galactic binaries, which dominate the confusion problem, the limited bandwidth of each individual signal means that the $`M^{ij}`$ is band diagonal, and thus easily inverted despite its large size. Since the search is now over the projected sub-space $`\{f_J,\theta _J,\varphi _J\}`$ of the full parameter space, the full Fisher matrix, $`\mathrm{\Gamma }_{ij}(\stackrel{}{x})`$, is replaced by the projected Fisher matrix, $`\gamma _{ij}(\stackrel{}{x})`$. The projection of the $`k^{\mathrm{th}}`$ parameter is given by $$\mathrm{\Gamma }_{ij}^{n1}=\mathrm{\Gamma }_{ij}^n\frac{\mathrm{\Gamma }_{ik}^n\mathrm{\Gamma }_{jk}^n}{\mathrm{\Gamma }_{kk}^n},$$ (29) where $`n`$ denotes the dimension of the projected matrix. Repeated application of the above projection yields $`\gamma _{ij}=\mathrm{\Gamma }_{ij}^{3N}`$. Inverting $`\gamma _{ij}`$ yields the same uncertainty estimates for the intrinsic parameters as one gets from the full Fisher matrix, but the covariances are much larger. The large covariances make it imperative that the proposal distributions use the eigenvalues and eigenvectors of $`\gamma _{ij}`$, as using the parameter directions themselves would lead to a slowly mixing chain. ## V F-Statistic MCMC We implemented an F-statistic based MCMC algorithm using the approach described in §III, but with the full likelihood replaced by the F-statistic and the full Fisher matrix replaced by the projected Fisher matrix. Applying the F-MCMC search to the same data set as before yields the results summarized in Figure 2 and Table 2. The recovered source parameters and signal-to-noise ratio ($`\mathrm{SNR}=10.4`$) are very similar to those found using the full 7-parameter search, but the F-MCMC estimates for the errors in the extrinsic parameters are very different. This is because the chain does not explore extrinsic parameters, but rather relies upon the F-statistic to find the extrinsic parameters that give the largest log likelihood based on the current values for the intrinsic parameters. The effect is very pronounced in the histograms shown in Figure 2. Similar results were found for other F-MCMC runs on the same source, and for F-MCMC runs with other sources. Typical burn-in times were of order 1000 iterations, and the proposal acceptance rate was around $`60\%`$. As expected, the F-MCMC algorithm gave shorter burn-in times than the full parameter MCMC, and a comparable mixing rate. It is interesting to compare the computational cost of the F-MCMC search to a traditional F-Statistic based search on a uniformly spaced template grid. To cover the parameter space of one source (which for the current example extends over the full sky and 100 frequency bins) with a minimal match ben1 of $`\mathrm{MM}=0.9`$ requires 39,000 templates neil\_ed . A typical F-MCMC run uses less than 1000 templates to cover the same search space. The comparison becomes even more lopsided if we consider simultaneous searches for multiple sources. A grid based simultaneous search for two sources using the F-statistic would take $`(39,000)^21.5\times 10^9`$ templates, while the basic F-MCMC algorithm typically converges on the two sources in just 2000 steps. As the number of sources in the model increases the computation cost of the grid based search grows geometrically while the cost of the F-MCMC search grows linearly. It is hard to imagine a scenario (other than quantum computers) where non-iterative grid based searches could play a role in LISA data analysis. While testing the F-MCMC algorithm on different sources we came across instances where the chain became stuck at secondary modes of the posterior. A good example occurred for a source with parameters $`(A,f,\theta ,\varphi ,\psi ,\iota ,\phi _0)`$=$`(1.4`$e-22$`,1.0020802\mathrm{mHz},`$ $`0.399,5.71,1.3,0.96,1.0)`$ and $`\mathrm{SNR}=16.09`$. Most MCMC runs returned good fits to the source parameters, with an average log likelihood of $`\mathrm{ln}=132`$, mean intrinsic parameter values $`(f,\theta ,\varphi )=(1.0020809\mathrm{mHz},0.391,5.75)`$ and $`\mathrm{SNR}=16.26`$. However, some runs locked into a secondary mode with average log likelihood $`\mathrm{ln}=100`$, mean intrinsic parameter values $`(f,\theta ,\varphi )=(1.0020858\mathrm{mHz},2.876,5.20)`$ and $`\mathrm{SNR}=14.15`$. It could sometimes take hundreds of thousands of iterations for the chain to discover the dominant mode. Figure 4 shows plots of the (inverted) likelihood $``$ and the log likelihood $`\mathrm{ln}`$ as a function of sky location for fixed $`f=1.0020802\mathrm{mHz}`$. The log likelihood plot reveals the problematic secondary mode near the south pole, while the likelihood plot shows just how small a target the dominant mode presents to the F-MCMC search. Similar problems with secondary modes were encountered in the $`f\varphi `$ plane, where the chain would get stuck a full bin away from the correct frequency. These problems with the basic F-MCMC algorithm motivated the embellishments described in the following section. ## VI Multiple Proposals and Heating The LISA data analysis problem belongs to a particularly challenging class of MCMC problems known as “mixture models.” As the name suggests, a mixture model contains a number of components, some or all of which may be of the same type. In our present study all the components are slowly evolving, circular binaries, and each component is described by the same set of seven parameters. There is nothing to stop two components in the search model from latching on to the same source, nor is there anything to stop one component in the search model from latching on to a blend of two overlapping sources. In the former instance the likelihood is little improved by using two components to model one source, so over time one of the components will tend to wander off in search of another source. In the latter instance it may prove impossible for any data analysis method to de-blend the sources (the marginal likelihood for the single component fit to the blended sources may exceed the marginal likelihood of the “correct” solution). The difficulties we encountered with the single source searches getting stuck at secondary modes of the posterior are exacerbated in the multi-source case. Source overlaps can create additional secondary modes that are not present in the non-overlapping case. We employed two techniques to speed burn-in and to reduce the chance of the chain getting stuck at a secondary mode: simulated annealing and multiple proposal distributions. Simulated annealing works by softening the likelihood function, making it easier for the chain to move between modes. The likelihood (8) can be thought of as a partition function $`Z=C\mathrm{exp}(\beta E)`$ with the “energy” of the system given by $`E=(sh\|sh)`$ and the “inverse temperature” equal to $`\beta =1/2`$. Our goal is to find the template $`h`$ that minimizes the energy of the system. Heating up the system by setting $`\beta <1/2`$ allows the Markov Chain to rapidly explore the likelihood surface. We used a standard power law cooling schedule: $$\beta =\{\begin{array}{cc}\beta _0\left(\frac{1}{2\beta _0}\right)^{t/T_c}\hfill & 0<t<T_c\hfill \\ \frac{1}{2}\hfill & tT_c\hfill \end{array}$$ (30) where $`t`$ is the number of steps in the chain, $`T_c`$ is the cooling time and $`\beta _0`$ is the initial inverse temperature. It took some trial and error to find good values of $`T_c`$ and $`\beta _0`$. If some of the sources have very high SNR it is a good idea to start at a high temperature $`\beta _01/50`$, but in most cases we found $`\beta _0=1/10`$ to be sufficient. The optimal choice for the cooling time depends on the number of sources and the initial temperature. We found that it was necessary to increase $`T_c`$ roughly linearly with the the number of sources and the initial temperature. Setting $`T_c=10^5`$ for a model with $`N=10`$ sources and an initial temperature of $`\beta _0=1/10`$ gave fairly reliable results, but it is always a good idea to allow longer cooling times if the computational resources are available. The portion of the chain generated during the annealing phase has to be discarded as the cooling introduces an arrow of time which necessarily violates the reversibility requirement of a Markov Chain. After cooling to $`\beta =1/2`$ the chain can explore the likelihood surface for the purpose of extracting parameter estimates and error estimates. Finally, we can extract maximum likelihood estimates by “super cooling” the chain to some very low temperature (we used $`\beta 10^4`$). The second ingredient in our advanced F-MCMC algorithm is a large variety of proposal distributions. We used the following types of proposal distribution: $`\mathrm{Uniform}(,\stackrel{}{x},i)`$ \- a uniform draw on all the parameters that describe source $`i`$, using the full parameter ranges, with all other sources held fixed; $`\mathrm{Normal}(,\stackrel{}{x})`$ \- a multivariate normal distribution with variance-covariance matrix given by $`3N\times \gamma (\stackrel{}{x})`$; $`\mathrm{Sky}(,\stackrel{}{x},i)`$ \- a uniform draw on the sky location for source $`i`$; $`\sigma `$-Uniform$`(,\stackrel{}{x},i)`$ \- a uniform draw on all the parameters that describe source $`i`$, using a parameter range given by some multiple of the standard deviations given by $`\gamma (\stackrel{}{x})`$. The $`\mathrm{Uniform}(,\stackrel{}{x},i)`$ and $`\mathrm{Normal}(,\stackrel{}{x})`$ proposal distributions are the same as those used in the basic F-MCMC algorithm. The $`\mathrm{Sky}(,\stackrel{}{x},i)`$ proposal proved to be very useful at getting the chain away from secondary modes like the one seen in Figure 4, while the $`\sigma `$-Uniform$`(,\stackrel{}{x},i)`$ proposal helped to move the chain from secondary modes in the $`f\varphi `$ or $`f\theta `$ planes. During the initial annealing phase the various proposal distributions were used in a cycle with one set of the bold distributions (Uniform, Sky and $`\sigma `$Uniform) for every 10 draws from the timid multivariate normal distribution. During the main MCMC run at $`\beta =1/2`$ the ratio of timid to bold proposals was increased by a factor of 10, and in the final super-cooling phase only the timid multivariate normal distribution was used. The current algorithm is intended to give a proof of principle, and is certainly far from optimal. Our choice of proposal mixtures was based on a few hundred runs using several different mixtures. There is little doubt that a better algorithm could be constructed that uses a larger variety of proposal distributions in a more optimal mixture. The improved F-MCMC algorithm was tested on a variety of simulated data sets that included up to 10 sources in a 100 bin snippet (once again we are using one year of observations). The algorithm performed very well, and was able to accurately recover all sources with $`\mathrm{SNR}>5`$ so long as the degree of source correlation was not too large. Generally the algorithm could de-blend sources that had correlation coefficients $`C_{12}=(h_1\|h_2)/\sqrt{(h_1\|h_1)(h_2\|h_2)}`$ below $`0.3`$. A full investigation of the de-blending of highly correlated sources is deferred to a subsequent study. For now we present one representative example from the 10 source searches. A set of 10 galactic sources was randomly selected from the frequency range $`f[0.999995,1.003164]`$ mHz and their signals were processed through a model of the LISA instrument response. The root spectral densities in the two independent LISA data channels are shown in Figure 5, and the source parameters are listed in Table 3. Note that one of the sources had a SNR below 5. The data was then search using our improved F-MCMC algorithm using a model with 10 sources (70 parameters). The annealing time was set at $`10^5`$ steps, and this was followed by a short MCMC run of $`2\times 10^4`$ steps and a super cooling phase that lasted $`2\times 10^4`$ steps. The main MCMC run was kept short as we were mostly interested in extracting maximum likelihood estimates. Figure 6 shows a trace plot of the chain that focuses on the frequencies of two of the model sources. During the early hot phase the chain moves all over parameter space, but as the system cools to $`\beta =1/2`$ the chain settles down and locks onto the sources. During the final super cooling phase the movement of the chain is exponentially damped as the model is trapped at a mode of shrinking width and increasing height. The list of recovered sources can be found in Table 3. The low SNR source ($`\mathrm{SNR}=4.9`$) was not recovered, but because the model was asked to find 10 sources it instead dug up a spurious source with $`\mathrm{SNR}=5.2`$. With two exceptions, the intrinsic parameters for the other 9 sources were recovered to within $`3\sigma `$ of the true parameters (using the Fisher matrix estimate of the parameter recovery errors). The two misses were the frequency of the source at $`f=1.00253`$ mHz (out by $`19\sigma `$) and the co-latitude of the the source at $`f=1.002632`$ mHz (out by $`6\sigma `$). It is no co-incidence that these misses occurred for the two most highly correlated sources ($`C_{9,10}=0.23`$). The full source cross-correlation matrix is listed in (VI). $`C_{ij}={\displaystyle \frac{(h_i\|h_j)}{\sqrt{(h_i\|h_i)(h_j\|h_j)}}}`$ $`=\left(\begin{array}{cccccccccc}\text{1}& \text{0.08}& \text{0}& \text{0.01}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}\\ \text{0.08}& \text{1}& \text{0.02}& \text{0.01}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}\\ \text{0}& \text{0.02}& \text{1}& \text{-0.06}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}\\ \text{0.01}& \text{0.01}& \text{-0.06}& \text{1}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{0}& \text{1}& \text{0}& \text{0}& \text{0.01}& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{1}& \text{-0.03}& \text{0.03}& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{-0.03}& \text{1}& \text{-0.05}& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{0}& \text{0.01}& \text{0.03}& \text{-0.05}& \text{1}& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{1}& \text{-0.23}\\ \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{-0.23}& \text{1}\end{array}\right)`$ (41) The MCMC derived maximum likelihood estimates for the the source parameters can be used to regress the sources from the data streams. Figure 7 compares the residual signal to the instrument noise. The total residual power is below the instrument noise level as some of the noise has been incorporated into the recovered signals. ## VII Model Selection In the preceding examples we used models that had the same number of components as there were sources in the data snippet. This luxury will not be available with the real LISA data. A realistic data analysis procedure will have to explore model space as well as parameter space. It is possible to generalize the MCMC approach to simultaneously explore both spaces by incorporating trans-dimensional moves in the proposal distributions. In other words, proposals that change the number of sources being used in the fit. One popular method for doing this is Reverse Jump MCMC rj , but there are other simpler methods that can be used. When trans-dimensional moves are built into the MCMC algorithm the odds ratio for the competing models is given by the fraction of the time that the chain spends exploring each model. While trans-dimensional searches provide an elegant solution to the model determination problem in principle, they can perform very poorly in practice as the chain is often reluctant to accept a trans-dimensional move. A simpler alternative is to compare the outputs of MCMC runs using models of fixed dimension. The odds ratio can then calculated using Bayes factors. Calculating the marginal likelihood of a model is generally very difficult as it involves an integral over all of parameter space: $$p_X(s)=p(s\|\stackrel{}{\lambda },X)p(\stackrel{}{\lambda },X)𝑑\stackrel{}{\lambda }.$$ (43) Unfortunately, this integrand is not weighted by the posterior distribution, so we cannot use the output of the MCMC algorithm to compute the integral. When the likelihood distribution has a single dominant mode, the integrand can be approximated using the Laplace approximation: $`p(\stackrel{}{\lambda },X)p(s\|\stackrel{}{\lambda },X)p(\stackrel{}{\lambda }_{\mathrm{ML}},X)p(s\|\stackrel{}{\lambda }_{\mathrm{ML}},X)`$ $`\times \mathrm{exp}\left({\displaystyle \frac{(\stackrel{}{\lambda }\stackrel{}{\lambda }_{\mathrm{ML}})F(\stackrel{}{\lambda }\stackrel{}{\lambda }_{\mathrm{ML}})}{2}}\right).`$ (44) where $`F`$ is given by the Hessian $$F_{ij}=\frac{^2\mathrm{ln}(p(\stackrel{}{\lambda },X)p(s\|\stackrel{}{\lambda },X))}{\lambda _i\lambda _j}\|_{\stackrel{}{\lambda }=\stackrel{}{\lambda }_{\mathrm{ML}}}.$$ (45) When the priors $`p(\stackrel{}{\lambda },X)`$ are uniform or at least slowly varying at maximum likelihood, $`F_{ij}`$ is equal to the Fisher matrix $`\mathrm{\Gamma }_{ij}`$. The integral is now straightforward and yields $$p_X(s)p(\stackrel{}{\lambda }_{\mathrm{ML}},X)p(s\|\stackrel{}{\lambda }_{\mathrm{ML}},X)\frac{(2\pi )^{D/2}}{\mathrm{det}F}.$$ (46) With uniform priors $`p(\stackrel{}{\lambda }_{\mathrm{ML}},X)`$=$`1/V`$, where $`V`$ is the volume of parameter space, and $`(2\pi )^{D/2}/\mathrm{det}F`$=$`\mathrm{\Delta }V`$, where $`\mathrm{\Delta }V`$ is the volume of the error ellipsoid. To illustrate how the Bayes factor can be used in model selection, we repeated the F-MCMC search described in the previous section, but this time using a model with 9 sources. The results of a typical run are presented in Table 4. The parameters of the 9 brightest sources were all recovered to within $`3\sigma `$ of the input values, save for the sky location of the source with frequency $`f=1.00253`$ mHz. It appears that confusion with the source at $`f=1.002632`$ mHz may have caused the chain to favour a secondary mode like the one seen in Figure 4. Using (46) to estimate the marginal likelihoods for the 9 and 10 parameter models we found $`\mathrm{ln}p_9(s)=384.3`$ and $`\mathrm{ln}p_{10}(s)=394.9`$, which gives an odds ratio of $`1:4\times 10^4`$ in favour of the 9 parameter model. In contrast, a naive comparison of log likelihoods, $`\mathrm{ln}_9=413.1`$ and $`\mathrm{ln}_{10}=425.7`$ would have favoured the 10 parameter model. It is also interesting to compare the output of the 10 source MCMC search to the maximum likelihood one gets by starting at the true source parameters then applying the super cooling procedure (in other words, cheat by starting in the neighborhood of the true solution). We found $`p_{\mathrm{cheat}}(s)=394.5`$, and $`\mathrm{ln}_{\mathrm{cheat}}=421.5`$, which tells us that the MCMC solution, while getting two of the source parameters wrong, provides an equally good fit to the data. In other words, there is no data analysis algorithm that can fully deblend the two highly overlapping sources. ## VIII Conclusion Our first pass at applying the MCMC method to LISA data analysis has shown the method to have considerable promise. The next step is to push the existing algorithm until it breaks. Simulations of the galactic background suggest that bright galactic sources reach a peak density of one source per five $`1/\mathrm{year}`$ frequency bins seth . We have shown that our current F-MCMC algorithm can handle a source density of one source per ten frequency bins across a one hundred bin snippet. We have yet to try larger numbers of sources as the current version of the algorithm employs the full $`D=7N`$ dimensional Fisher matrix in many of the updates, which leads to a large computational overhead. We are in the process of modifying the algorithm so that sources are first grouped into blocks that have strong overlap. Each block is effectively independent of the others. This allows each block to be updated separately, while still taking care of any strongly correlated parameters that might impede mixing of the chain. We have already seen some evidence that high local source densities pose a challenge to the current algorithm. The lesson so far has been that adding new, specially tailored proposal distributions to the mix helps to keep the chain from sticking at secondary modes of the posterior (it takes a cocktail to solve the cocktail party problem). On the other hand, we have also seen evidence of strong multi-modality whereby the secondary modes have likelihoods within a few percent of the global maximum. In those cases the chain tends to jump back and forth between modes before being forced into a decision by the super-cooling process that follows the main MCMC run. Indeed, we may already be pushing the limits of what is possible using any data analysis method. For example, the 10 source search used a model with 70 parameters to fit 400 pieces of data (2 channels $`\times `$ 2 Fourier components $`\times `$ 100 bins). One of our goals is to better understand the theoretical limits of what can be achieved so that we know when to stop trying to improve the algorithm! It would be interesting to compare the performance of the different methods that have been proposed to solve the LISA cocktail party problem. Do iterative methods like gCLEAN and Slice & Dice or global maximization methods like Maximum Entropy have different strengths and weakness compared to MCMC methods, or do they all fail in the same way as they approach the confusion limit? It may well be that methods that perform better with idealized, stationary, Gaussian instrument noise will not prove to be the best when faced with real instrumental noise. ###### Acknowledgements. This work was supported by NASA Cooperative Agreement NCC5-579.
warning/0506/hep-th0506046.html	ar5iv	text	# 1 Initial Considerations ## 1 Initial Considerations The similarity of expressions from $`\text{SLE}(\kappa ,\rho )`$ with correlators in the Coulomb gas formalism was (probably) first noticed by S. Chakravarty . His questions for an explanation of this “coincidence” led to some (almost) unpublished notes . An independent explanation of $`\text{SLE}(\kappa ,\rho )`$ from the CFT perspective was given in . In this paper we will show, thereby relying on the results in , that $`\text{SLE}(\kappa ,\rho )`$ arises naturally when one considers random growing compacts in polygons, and how this fits into the context of conformal field theories. We will also outline, that there are more meaningful stochastic processes of SLE($`\kappa ,\rho `$) type, that can be derived from physical considerations. For the mathematical details concerning the “fluctuating polygons”, see . For the relations of SLE to diffusion processes on moduli spaces and / or general CFT, see . ### 1.1 SLE($`\kappa ,\rho `$) $`\text{SLE}(\kappa ,\rho )`$ was introduced on mathematical grounds, as a generalisation of “ordinary” SLE, in the landmark paper by Lawler, Schramm and Werner, and further studied in . (Extensive mathematical details on SLE can be found, e.g. in ). Stochastic Loewner evolution (or SLE) as introduced by Schramm in describes random growing compacts, in simply connected planar domains, which correspond (supposedly) to the conformally invariant scaling limit of discrete random simple curves that also satisfy a Markovian-type property. Then by the two above properties (plus a reflection symmetry) SLE is canonical in the sense that there exists only a one-parameter family of random non-selfcrossing curves $`\gamma `$ with these properties. The dynamical way to describe the measures, is by solving Löwner’s equation with a driving function given in terms of Brownian motion. So for the upper half-plane $``$, and $`\kappa 0`$, consider for each $`z\overline{}`$ the ordinary differential equation $$_tg_t(z)=\frac{2}{g_t(z)W_t},g_0(z)=z,$$ (1) where $`W_t=\sqrt{\kappa }B_t`$, and $`B_t`$ is a one-dimensional standard Brownian motion. Let $`T_z`$ be the duration for which this equation is well defined, i.e. $`T_z=sup\{t:inf_{s[0,t]}\|g_t(z)W_t\|>0\}`$, and set $`K_t=\{z:T_zt\}`$. Then one can show that $`g_t`$ is a conformal map from $`\backslash K_t`$ onto $``$ with $`lim_z\mathrm{}(g(z)z)=0`$. It can also be shown that with probability one the random growing compact set $`K_t`$ is generated by a random non-selfcrossing curve $`t\gamma _t`$ in the sense that $`\backslash K_t`$ is the unbounded component of $`\backslash \gamma [0,t]`$. $`\gamma `$ is a random curve connecting the boundary points $`0`$ and $`\mathrm{}`$ and is called chordal $`\text{SLE}_\kappa `$ in $``$ from $`0`$ to $`\mathrm{}`$. For calculations involving SLE conformal invariance is a powerfull tool as it is always permissible to choose the geometrically most convenient configuration to do a given calculation, where the solution depends only on the conformal equivalence class, or the moduli, of the configuration. Now, let $`z_1<z_2<\mathrm{}<z_n`$ be real numbers, all distinct from $`0`$. Consider the system of stochastic differential equations $`dW_t`$ $`=\sqrt{\kappa }dB_t+{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{\rho _k}{W_tZ_t^k}}dt`$ $`dZ_t^k`$ $`={\displaystyle \frac{2}{Z_t^kW_t}}dt,k=1,\mathrm{},n,`$ (2) with $`W_0=0,Z_0^1=z_1,\mathrm{},Z_0^n=z_n`$, and where $`B_t`$ is a one-dimensional standard Brownian motion. Then, at least up to some small time $`t`$, the solution will exist. As above, let $`g_t(z)`$ be the solution to (1). Then the family of conformal maps $`g_t`$ is called $`\text{SLE}(\kappa ,\rho )`$ in the upper half-plane from $`(0,z_1,\mathrm{},z_n)`$ to $`\mathrm{}`$. ### 1.2 From Physics Let us start with aspects of Liouville field theory , (for an early review ), that is related to the problem of quantum gravity. There, one of the major tasks is to properly integrate over all metrics modulo diffeomorphisms. So, let us consider a two-dimensional surface $`M`$, possibly with boundary $`M`$ and a Riemannian metric $`g`$, i.e. a bordered Riemann surface. In the conformal gauge any metric can be written as $$g=e^{\gamma \varphi }g_0,$$ (3) where $`\gamma `$ is a parameter and $`g_0`$ is the “reference metric”, which also determines a point $`[g_0]`$ in the moduli space. The field $`\varphi `$ is known as the Liouville mode. For a closed surface $`M`$ the underlying field theory is given by a bulk action $$S[g_0,\varphi ]:=\frac{1}{8\pi }_M\sqrt{g_0}\left((\varphi )^2+Q\varphi R(g_0)+\frac{\mu }{\gamma ^2}e^{\gamma \varphi }\right)d^2z,$$ (4) with the coupling constant $`\gamma (\mathrm{}=\gamma )`$, $`R()`$ the scalar curvature of $`(M,g_0)`$ (in dimension two $`R`$ is twice the sectional curvature) and the cosmological constant $`\mu >0`$ . Let us point out, that in the above action (4), the $`Q`$-term represents the action in the Coulomb gas formalism (CGF) on closed Riemann surfaces, i.e. a Gaussian conformal field theory in the presence of an imaginary background charge with $`Q=2i\alpha _0`$. Its presence leads to a modified stress-energy tensor $$T=\frac{1}{2}\varphi \varphi +i\alpha _0^2\varphi ,$$ (5) which generates a Virasoro algebra with central charge $$c=112\alpha _0^2.$$ Further, let us mention that (4) also represents a quantum conformal field theory. But contrary to “standard” CFT where we have the state-field correspondence, this does not hold any longer in Liouville field theory. Here the primary operators are of the form $`e^{\alpha \varphi }`$ with conformal weight $`\frac{1}{2}(\alpha \frac{Q}{2})^2+\frac{Q^2}{8}`$, and the set of operators and the set of states are distinct. Now let us come to the case of surfaces with boundaries. There are two major new aspects. First the bulk action has to be extended by contributions from the boundary and second, one has to fix boundary conditions. So, let us define the boundary action as $$S_LS_{\text{bulk}}+\frac{Q}{8\pi }_M\varphi k\|dz\|+\frac{\lambda }{4\pi \gamma ^2}_Me^{\frac{1}{2}\gamma \varphi }\|dz\|,$$ (6) where $`k`$ is the geodesic curvature of the boundary, $`\|dz\|`$ the line element and $`\lambda `$ the boundary cosmological constant. To have a well posed variational problem, the possible boundary conditions are Dirichlet or Neumann, or combinations of the two. Then the bulk equation of motion $`\frac{\delta S}{\delta \varphi }`$ is: $`R(g)`$ $`=`$ $`{\displaystyle \frac{\mu }{2}},\text{which is equivalent to}`$ (7) $`\mathrm{\Delta }\gamma \varphi `$ $`=`$ $`{\displaystyle \frac{\mu }{2}}e^{\gamma \varphi }+R(g_0),`$ i.e. the metric $`g=e^{\gamma \varphi }g_0`$ has constant negative curvature. The stress-energy tensor is found by varying the action with respect to the reference metric, i.e. $`T_{ab}=2\pi \frac{\delta S}{\delta g_0^{ab}}`$, which results in $`T_{z\overline{z}}`$ $`=`$ $`0,`$ (8) $`T_{zz}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\varphi )^2+{\displaystyle \frac{1}{2}}Q^2\varphi .`$ As $`\varphi `$ is a component of a metric as well, it transforms under conformal mappings $`zw=f(z)`$ like $$\varphi \varphi +\frac{1}{\gamma }\mathrm{log}\left\|\frac{dw}{dz}\right\|^2.$$ (9) In particular, the $`U(1)`$ current $`_z\varphi `$ transforms as $$_z\varphi \frac{dw}{dz}_w\varphi +\frac{d}{dz}\frac{1}{\gamma }\mathrm{log}\left\|\frac{dw}{dz}\right\|,$$ (10) and the stress tensor $`T_{zz}`$ as $$T_{zz}\left(\frac{dw}{dz}\right)^2T_{ww}+\frac{1}{\gamma ^2}\{w;z\}.$$ (11) Here, $`\{w;z\}`$ denotes the Schwarzian derivative. The equations of motion in the bordered case, with ($`\delta \varphi \|M=0`$) are: $$\frac{(\gamma \varphi )}{𝐧}+k+\frac{\lambda }{2}e^{\frac{1}{2}\gamma \varphi }=0,$$ (12) where the n-term denotes the normal derivative. We note, that like in the case of the bulk, there are also vertex operators on the boundary with some conformal weight $`\mathrm{\Delta }_M=2\alpha ^2+Q\alpha `$. In the above discussions we tacitly assumed, that the metric has no singularities neither in the bulk nor on the boundary. We shall shortly see, how that changes things. ### 1.3 The Uniformisation Problem The solutions of the classical equations of motion (7) and (12), i.e. the Liouville equation, are intrinsically related to the uniformisation problem of Riemann surfaces. It states, that every Riemann surface is conformally equivalent to either the Riemann sphere, the upper half-plane $``$ or to a quotient of $``$ by some discrete subgroup $`\mathrm{\Gamma }SL(2,)`$. The “fundamental” solution of the Liouville equation for $``$ is the Poincaré metric with constant negative curvature $`1`$. The classical solutions in Euclidean space of the Liouville equations are of the form $$e^{\gamma \varphi }\|dz\|^2=\frac{4}{\mu }\frac{A\overline{}B}{\left(A(z)B(\overline{z})\right)^2}\|dz\|^2$$ (13) with $`A`$ and $`B`$ some (locally) defined functions of $`z,\overline{z}`$. If $`X=/\mathrm{\Gamma }`$ denotes the quotient by a discrete subgroup, then there exists a natural projection $`\pi :X`$ with an “inverse” map $$f:X,$$ (14) depending on the moduli of $`X`$. In terms of the inverse map $`f`$ the solution of the field equation (13) has energy-momentum tensor (Fuchsian projective connection) satisfying: $$T_{zz}=\frac{1}{\gamma ^2}\{f;z\}.$$ (15) So-far we have excluded metric singularities, whose presence either in the interior or on the boundary, we will now permit. Then, depending on the conjugacy classes of the monodromy of $`A`$ and $`B`$ and of the nature of the metric singularity, there are three classes of local solutions: elliptic, parabolic and hyperbolic. For the present paper the elliptic case is the important one. 1. Elliptic : the solution has a curvature singularity, and so (e.g. here with a curvature source at $`z=0`$ and for $`a`$) the Liouville equation reads $$\frac{1}{4\pi }\mathrm{\Delta }\varphi \frac{\mu }{8\pi \gamma }e^{\gamma \varphi }+\frac{1a}{\gamma }\delta ^{(2)}(z)=0,$$ (16) Geometrically this is the situation corresponding to conical singularities / corners (orbifolds). 2. Parabolic : corresponds to punctured Riemann surfaces / surfaces with infinite cusps. 3. Hyperbolic : corresponds to a constant negative curvature metric on the annulus, i.e. “plumbing fixture metric”. There is a simple way of producing conical singularities, more precisely corners. Let us consider some polygon (details will follow in later parts of the paper). Then by the inverse of the standard Schwarz-Christoffel mapping we can biholomorphically map the polygon onto the upper half-plane, such that the vertices get mapped onto the real axis. By pulling-back the reference metric on the polygon, we get a metric with corner singularities on $``$. ### 1.4 Conical singularities and the space of Polygons Let us consider a bordered Riemann surface $`M`$. A (real) divisor on $`M`$ is the formal sum $$𝜷=\underset{i}{}\beta _ip_i$$ where the $`p_iM`$ are points and $`\beta _i`$. The discrete set $`\{p_i\}`$ is the support of $`𝜷`$ and the number $`\|𝜷\|:=_i\beta _i`$ is the degree of the divisor. We shall have the following conditions on the divisor: $$\beta _i>1\text{if}p_iM\text{and}\beta _i>\frac{1}{2}\text{if}p_iM.$$ (17) We shall call a simply connected domain $`D`$ with a divisor $`\{(p_1,\beta _1),\mathrm{},(p_n,\beta _n)\}`$, a weighted domain. Then a conformal metric $`ds^2`$ on $`M`$ represents the divisor $`𝜷`$ if $`ds^2`$ is a $`C^2`$-Riemannian metric on $`M\mathrm{supp}(𝜷)`$ such that if $`z_i`$ is a local coordinate on a neighbourhood $`U_i`$ of $`p_i`$, then there exists a continuous function $`u:U_i`$, of class $`C^2`$ on $`U_i\{p_i\}`$, such that on $`U_i`$: $$\{\begin{array}{cc}ds^2=e^{2u}\|z_ia_i\|^{2\beta _i}\|dz_i\|^2\hfill & \text{if}p_iM,\hfill \\ ds^2=e^{2u}\|z_ia_i\|^{4\beta _i}\|dz_i\|^2\hfill & \text{if}p_iM,\hfill \end{array}$$ (18) where $`a_i=z_i(p_i)`$. The point $`p_i`$ is called a conical singularity of angle $`\theta _i=2\pi (\beta _1+1)`$ if $`p_iM`$ and a corner of angle $`\phi _i=2\pi (\beta _i+\frac{1}{2})`$ or of exterior angle $`2\pi \beta _i`$ if $`p_iM`$. In both cases, we shall say that $`ds^2`$ has a singularity of order $`\beta _i`$ at $`p_i`$. Riemann surfaces with conical singularities are generalised Riemann surfaces (GRS) but the natural morphisms for GRS’s are still conformal mappings, which topologically are covering maps in the sense of 2 dimensional orbifold theory. The problem of prescribing curvature for such surfaces, has been studied in . Now, the Gauss-Bonnet formula does still hold in this situation, and it states for a GRS with divisor $`(M,𝜷)`$ that: $$\frac{1}{2\pi }_MR\text{dvol}+\frac{1}{2\pi }_Mk\|dz\|=\chi (M,𝜷),$$ (19) where the Euler characteristic of $`(M,𝜷)`$ is defined by $`\chi (M,𝜷)\chi (M)+\|𝜷\|`$ with $`\chi (M)`$ the topological Euler characteristic of $`M`$. In a physical context this corresponds to some conservation law, e.g. charge conservation. As the above discussion shows, polygons are natural objects to consider, if one wants to deal with conical singularities (corners). Therefore let us briefly mention some facts about the set $`𝒫_n`$ of all polygons with $`n3`$ distinguished vertices in the complex plane $``$, whose sides have non-negative length. We shall allow for all possible degenerations of the polygons, with the exception of the degeneration to a single point. On the complex plane $``$ we shall consider the usual Euclidean metric $`\|dz\|^2`$. Two polygons are called equivalent if there is an orientation preserving similarity of the complex plane, which maps vertices of one polygon to those of the other one. We know, that these conformal mappings are given by a global linear transformation of the form: $$f(z):=az+b,\text{where}(a,b)^{}.$$ It is also well known that the infinitesimal generators of these transformations are the differential operators: $`_z`$ (translations) and $`z_z`$ (dilatations and rotations) If we denote the edges of the $`n`$-gon $`D`$ by $`e_1,\mathrm{},e_n`$ and its vertices by $`v_1,\mathrm{},v_n`$, then we have the basic relation from vector calculus $`\stackrel{}{e}_j=v_{j+1}v_j`$. So the space $`𝒫_n`$ is canonically isomorphic to the complex projective space $`^{n2}`$. To see this, just consider the hyperplane $`H^n`$, defined by: $$H:=\{(e_1,\mathrm{},e_n)^n:e_1+\mathrm{}+e_n=0\}.$$ Therefore, the space $`𝒫_n`$ is connected and is naturally endowed with the Fubini-Study metric. Hence, every polygon $`D𝒫_n`$ can be continuously deformed to any polygon $`P𝒫_n`$. ### 1.5 The Correlator Toolbox As the present article deals with stochastic processes of $`\text{SLE}(\kappa ,\rho )`$ type, which are basically stochastic multi-particle systems, where the individuals are confined to the boundary of some surface, originally the real line, in a physical context we have to consider, what kind of particles we are dealing with, e.g. what quantum numbers they posses. Whereas in probability theory one might just see a random dynamical system, the physicist perceives them (in the SLE context) rather as objects (quanta) belonging to an underlying field theory. As our previous discussion showed, there are several quantum field theories which are related to the realm of conformal symmetry. They are derived from some classical action and quantized via the path-integral formalism, which in turn, and very loosely speaking, corresponds to the Itô Integral. Below, the situation is depicted schematically, where $`\varphi `$ stands for the relevant fields (or just a dummy variable): $$\underset{\text{BLFT}}{\underset{}{\underset{\text{CGF}}{\underset{}{\underset{\text{BCFT}}{\underset{}{(\varphi )^2}}QR\varphi }}e^{\gamma \varphi }}}\text{+ boundary terms + boundary conditions}+\text{ghosts}$$ Further, we also saw, that the quantum field theories discussed, have various characteristic operators of type, e.g.: * bulk / boundary vertex operators * boundary condition changing operators (Cardy type) * twist fields: Dirichlet$``$Neumann (“dual resonance theory” type) but also * bulk and boundary curvature sources of elliptic or parabolic type As the fundamental quantities in a quantum field theory are given by correlators, e.g. $`Z`$ the partition function, we could try to build out of the above operators an arbitrary “correlator burger”, $$\mathrm{}\text{plug in operators}\mathrm{}$$ In the simplest case, it could be derived from the situation in Fig. 2, where we have on the real axis different boundary conditions, Dirichlet and Neumann or discontinuously changing boundary conditions of the same type, curvature sources in the bulk or on the boundary. But, as it is known, conformal field theories are controlled by tight algebraic structures, e.g. the central charge and the dimensions of the fields are dictated by representation theory of, e.g. the Virasoro algebra. Therefore, not all of the possible insertions would lead to well defined expressions. Further, as it was already known in the early days of CFT, that when a conformal field theory is considered on a Riemann surface, the moduli will enter explicitly . Examples, of “natural” Riemann surfaces with non-trivial moduli are the unit disc with $`10^{26}`$ marked points, the annulus or an $`n`$-connected domain, to mention some of the planar ones. So, as an illustration, for the Ising model in highest-weight representation with $`c=1/2`$ there exists a degeneracy at level two with highest-weight state $`h=1/2`$, i.e. $$[L_2\frac{3}{4}(L_1)^2]\|h=1/2=0,$$ (20) and therefore this state must decouple from the other states in an unitary irreducible representation. This leads, as a particularity of CFT, to a differential equation for correlation functions, i.e. there is a correspondence of null states in the representation space and linear differential operators whose order is given by the grade of the null state. In the case of a non-trivial Riemann surface, the Ward identities acquire an additional term, which describes the dependence of correlation functions on the moduli of the underlying surface (as a function of the topological invariants, e.g. genus, boundary components, marked points). Although the local algebraic structure (20) remains the same on the Riemann surfaces, the identity (20) becomes an operator identity, $$(L_2(w)\frac{3}{4}(L_1(w))^2)\varphi _{h=1/2}(w)=0.$$ (21) where $`\varphi `$ is a conformal field of dimension $`h`$. Now, in the SLE / CFT context, we are interested in certain choices of operators within a correlator, that geometrically seen, create at least one simple curve, e.g. a domain wall, and in possible conformally invariant probability measures supported by these curves. Of course, the measures will depend also on the presence of the other fields. At this point one may consult . By the Loewner mapping, i.e. by cutting the surface, we can evolve the correlator, to obtain the relevant driving process, and hence by pull-back the probability measure itself. However the stochastic driving process now lives on the appropriate moduli space, and is not a simple one-dimensional Brownian motion, any longer. In case, of conformal invariance and a Markovian-type property the process will be a Markov process, see . A look at Fig. (2), immediately reveals that the situation can be naturally extended to the case, where we have also insertions of some bulk fields and / or higher loop diagrams. Therefore, a natural general correlator, which is still describable within the generalised SLE framework, leads to processes of the form $$\text{SLE}(\kappa ,\stackrel{}{\rho }_{\text{boundary}},\stackrel{}{\rho }_{\text{bulk}}).$$ (22) The consideration of these, should give new and interesting classes of stochastic processes with associated measures. ## 2 An Illustration: Diffusing Polygons Let us now illustrate the preceding discussion with an example. To do so, let us consider a critical statistical mechanics model defined on a polygon $`D`$, with changing boundary conditions at points $`A`$ and $`B`$, see Fig. (3). As in the usual SLE set-up, this generates a domain wall, that connects the points $`A`$ and $`B`$ and in principle is described by a “sort of chordal SLE”. We will now derive the driving markov process in this situation from purely geometric considerations. ### 2.1 Schwarz-Christoffel Formula. Now let the consecutive vertices be $`p_1,\mathrm{}p_n`$ in positive cyclic order. The angle at $`p_k`$ is $`\alpha _k\pi `$, $`0<\alpha _k<2`$, and the outer angle is $`\beta _k\pi =(1\alpha _k)\pi `$, $`1<\beta _k<1`$. We note that $$\beta _1+\mathrm{}+\beta _n=2,$$ (23) and that the polygon is convex if and only if all $`\beta _k>0`$. We will call the pairs $`(p_k,\beta _k)`$ the corners of the polygon. Further, let $`f`$ be a conformal map from $`D`$ onto the upper half-plane $``$, with $`z_k=f(p_k)`$ and such that none of the $`z_k`$ equals $`\mathrm{}`$. Then for $`z`$ define the Schwarz-Christoffel mapping $$SC(z)=SC\left[\begin{array}{cc}z_1,\mathrm{},z_n& z\\ \beta _1,\mathrm{},\beta _n& z^{}\end{array}\right]=_z^{}^z\underset{k=1}{\overset{n}{}}(zz_k)^{\beta _k}dz,$$ (24) where the powers $`(zz_k)^{\beta _k}`$ denote analytic branches in $``$. Note that $$SC^{}(z)=\underset{k=1}{\overset{n}{}}(zz_k)^{\beta _k},$$ (25) and $$\frac{SC^{\prime \prime }(z)}{SC^{}(z)}=\underset{k=1}{\overset{n}{}}\frac{\beta _k}{zz_k}.$$ (26) Then it is well known that for some constants $`a,b`$, $`f^1=aSC+b`$, see . This result extends to the case when the polygon is allowed to have slits, i.e $`\beta _k=1`$ for some $`k`$. Slits are counted as double edges of the boundary polygon, traversed in positive cyclic order. A vertex, when considered as a boundary point, may then occur multiple times, corresponding to different prime ends. Henceforth, a corner $`(p_k,\beta _k)`$ will always be a pair consisting of a prime end $`p_k`$ located at a vertex together with the exterior angle $`\beta _k`$ associated to the prime end $`p_k`$. The formula (24) then remains unchanged if $`z_k=f(p_k)`$. In case $`f(p_k)=\mathrm{}`$ for one $`k`$, then (24) needs to be adjusted by simply dropping the factor with exponent $`\beta _k`$. Finally, formula (24) continues to hold if $`D`$ is unbounded or one (or several) corners are at $`\mathrm{}`$, provided that the angles at $`\mathrm{}`$ are appropriately defined, see . The interior angle at $`\mathrm{}`$ is chosen in $`[2\pi ,0]`$. So, polygon will refer to these “generalised” polygons. In the case of polygons, we still have, that if $`D`$ is a polygon with corners $`(p_1,\beta _1)`$ to $`(p_n,\beta _n)`$ in positive cyclical order, and $`\gamma `$ is a Jordan arc contained in $`D`$ except for one endpoint which lies on the interior of a side $`S`$ of $`D`$, then there is a conformal map $`f`$ from $`D\backslash \gamma `$ onto a polygon $`D^{}`$ such that $`(f(p_1),\beta _1),\mathrm{},(f(p_n),\beta _n)`$ are the corners of $`D^{}`$ in positive cyclical order. If we require $`f(S\gamma )[0,1]`$, then $`f`$ is unique. ## 3 SLE$`(\kappa ,\rho )`$ and Polygon motion Let $`(p_1,\beta _1),\mathrm{},(p_n,\beta _n)`$ be the corners of a polygon $`D`$. Denote $`f`$ a conformal map from $`D`$ onto the upper half-plane and set $`z_k=f(p_k)`$, $`1kn`$. We assume the points $`z_k`$ are all finite and distinct from $`0`$. For $`\kappa >0`$ set $$\rho _k=\frac{\kappa }{2}\beta _k,k=1,\mathrm{},n.$$ (27) In particular, $`\kappa /2\rho _k\kappa /2`$. Suppose that $`(W_t,Z_t^1,\mathrm{}Z_t^n)`$ is a solution to (1.1). For $`z`$ in the upper half-plane, set $$SC_t(z)=SC\left[\begin{array}{cc}Z_t^1,\mathrm{},Z_t^n& z\\ \beta _1,\mathrm{},\beta _n& 0\end{array}\right].$$ Then $`zSC_t(z)`$ extends continuously to the real axis with the points $`Z_t^k`$ removed and is differentiable there as a function of $`t`$. In particular, if $`W_sZ_s^1,\mathrm{},Z_s^n`$ for $`s[0,t]`$, then we may define $$f_t(z)=SC_t(z)_0^t(_sSC_s)(W_s)𝑑s.$$ (28) Note that $`f_t`$ maps $``$ onto a polygon while the function $`f`$ as previously introduced or in Fig. (3), maps a polygon onto the upper half-plane, i.e. is the inverse Schwarz-Christoffel mapping. Define the stopping time $`\sigma `$ by $$\sigma =sup\{t:W_s,Z_s^1,\mathrm{},Z_s^n\text{are all distinct for }0st\}.$$ Then the process $`U_tf_t(W_t)`$ is a martingale for $`t<\sigma `$ and if $$A_t\kappa _0^t\left(SC_s^{}(W_s)\right)^2𝑑s$$ and $`\tau (t)`$ is defined by $`A_{\tau (t)}=t`$, then $`tU_{\tau (t)}`$ is a standard Brownian motion. We further note, that the motion of the corners of the polygon $`f_t()`$ is differentiable. It is important to point out, that if we begin with an arbitrary SLE($`\kappa ,\rho `$), i.e. with a choice of points $`z_1,\mathrm{},z_2`$ and weights $`\rho _1,\mathrm{},\rho _n`$, then the results of this section will still hold. However, in this case the Schwarz-Christoffel mapping $`SC`$ is no longer to be schlicht, but it still maps the intervals $`[z_k,z_{k+1}]`$ onto straight line segments. But, by considering the associated Riemann surface to the analytic function $`SC`$, we can still interpret the image $`SC()`$ as a Riemannian domain. So, in the case of SLE($`2,(1,1))`$, up to a normalisation, this corresponds to the map $`z^33z`$, which can be understood in terms of a $`3`$-fold cover, cf. Set now $`D_t=f_t()`$, and denote $$(q,u)D_t\times D_tk_{D_t}(q,u)$$ the Poisson kernel of $`D_t`$. If $`pD_t`$, denote $`_2H_{D_t,p}(q,u)`$ the analytic function in $`q`$ whose real part is $`_2k_{D_t}(q,u)`$ and which satisfies $$\underset{qp}{lim}_2H_{D_t,p}(q,u)=0.$$ Then the main statement is ###### Theorem 3.1 (Loewner evolution in polygons). Denote $`K_t`$ the hull of an $`\text{SLE}_\kappa (\rho )`$ in the upper half-plane and $`g_t:\backslash K_t`$ the normalised uniformising map. Then $`h_tf_tg_tf_0^1:D\backslash f_0(K_t)D_t`$ satisfies $$_t\mathrm{ln}h_t^{}(z)=f_t(W_t)^2_2H_{D_t,f_t(\mathrm{})}(h_t(z),f_t(W_t)).$$ (29) ## 4 SLE coupled to Gravity As it is know, a stochastic processes is linked to some operator, which intrinsically depends upon a metric. The operator which determines what Brownian motion should be, is the Laplace-Beltrami operator. In local coordinates, applied to a function $`\phi `$, it reads as $$\mathrm{\Delta }\phi =\frac{1}{\sqrt{g}}\frac{}{x^j}\left(\sqrt{g}g^{ij}\frac{\phi }{x^i}\right).$$ So, instead of mapping $`\text{SLE}(\kappa ,\rho )`$ into polygons we can also stay in the upper half-plane and couple it to a fluctuating background metric. Indeed, $`f_t:D_t`$ is an immersion. If we endow $`D_t`$ with the Euclidean metric, then the pull-back metric via $`f_t`$ on $``$ is $$g_{ij}=\delta _{ij}\|f_t^{}(z)\|^2,i,j=1,2,$$ where the indices 1 and 2 refer to the real and imaginary coordinate, respectively. If $`\mathrm{\Gamma }=(\mathrm{\Gamma }_{jk}^i)`$ denotes the Levi-Civita connection for this metric, then the (2-dimensional) Brownian motion $`\stackrel{~}{W}`$ for the metric $`(g_{ij})`$ solves the stochastic differential equation $$d\stackrel{~}{W}_s^i=\sigma _j^i(\stackrel{~}{W}_s)dB_s^j\frac{1}{2}g^{kl}(\stackrel{~}{W}_s)\mathrm{\Gamma }_{kl}^i(\stackrel{~}{W}_s)ds,$$ see . Here $`g^1=(g^{kl})`$ is the inverse coefficient matrix of $`g`$ and $`\sigma `$ is a square root of $`g^1`$ (i.e. $`\sigma \sigma ^T=g^1`$), and we apply the Einstein summation convention. For our particular metric $`g`$ we find $$\mathrm{\Gamma }_{11}^1=\mathrm{\Gamma }_{22}^1=\mathrm{}\left(\frac{f_t^{\prime \prime }}{f_t^{}}\right),$$ see . The boundary $`=`$ is a one-dimensional submanifold of $`\overline{}`$. The metric $`g`$ on $``$ thus induces the metric $`(f_t^{}(x))^2dx^2`$ on $``$. A (one-dimensional) Brownian motion $`W`$ relative to this metric solves the stochastic differential equation $$dW_s=\frac{dB_s}{f_t^{}(W_s)}\frac{1}{2(f_t^{}(W_s))^2}\underset{j=1}{\overset{n}{}}\frac{\beta _j}{W_sZ_t^k}ds.$$ (30) We now couple the metric to the Brownian motion $`W`$ via $$dZ_t^k=\frac{2}{\kappa (f_t^{}(W_t))^2(Z_t^kW_t)}dt,k=1,\mathrm{},n.$$ (31) Then, after a timechange, (30) and (31) become the $`\text{SLE}(\kappa ,\rho )`$-system (1.1) with the convention $`\rho _j=\kappa \beta _j/2`$. #### Acknowledgements R. F. would like to thank Shoibal Chakravarty for the questions asked and the discussions. Gastón Giribert he thanks for helpful explanations. John Cardy he thanks for discussing the notes and for general discussions. The research of the first author was supported by NSA grant H98230-04-1-0039. The research of the second author was supported by a grant of the Max-Planck-Gesellschaft.
warning/0506/quant-ph0506081.html	ar5iv	text	# 1 Introduction ## 1 Introduction Quantum computing is intensive developing intersection of physics and informatics since (Feynman 1982; Shor 1994; Grover 1997; Preskill 1998). But schemes of universal quantum computers working by qubits meet formidable difficulties in their realization as real physical devices (Dyakonov 2003). In this note we develop another approach based on analogue realization of quantum computing. We concentrate our attention on some particular NP-hard problems. For a given problem we can try to construct a quantum analogue machine that resolves this problem by means of time evolution. Recall (Garey and Johnson 1979) that, roughly speaking, a problem with instances $`I`$ lies in the class NP if there is a polynomial time $`P(\|I\|)`$ algorithm checking a solution (if this solution is provided) where $`\|I\|`$ is the size of the input $`I`$. The problem is NP-hard if $`P=NP`$ whenever this problem belongs to $`P`$. The list of NP-complete problems contains many important problems of number theory, graph theory, logic etc. Some of the most known NP-complete problems are travelling salesman problem, satisfiability, knapsack, matrix cover etc. (see a list in (Garey and Johnson 1979)). Let us observe that all the NP-complete problems, are, in a sense, polynomially equivalent. This means that there are polynomial time reductions between different problems. An analogue quantum machine or resolving of combinatorial search NP-hard problems has been proposed by (Farhi et al. 2000). Let us consider 3-SAT problem in which we deal with a family of $`N`$ clauses of the form $`(B_{i_1}B_{i_2}\neg B_{i_3})`$ or $`(B_{i_1}\neg B_{i_2}\neg B_{i_3})`$ or $`(B_{i_1}B_{i_2}B_{i_3})`$ etc., where $`B_j`$ are boolean variables taking values True or False. The problem is whether one can assign such values to all the variables that all the clauses will become true. The approach of (Farhi et al. 2000) uses quantum adiabatic theorem. The authors describe a formal Hamiltonian such that its ground state gives us a solution of 3-SAT problem. The bound of the actual speed-up yielded by this algorithm depends on the spectral gap between the ground state and the first excited state. If this gap is exponentially small the time of solving is exponentially large. It is a very difficult problem to estimate the value of this gap and thus it is not easy to show, in general, that such an algorithm really gives a quantum acceleration. (Another NP-complete problem that can be naturally associated with a Hamiltonian is the problem of minimization of energy of spin glass with a large number of spins (see (Garey and Johnson 1979, p. 282.)) The second difficulty with such approach is that it is not obvious to realize physically qubits and Hamiltonian with very nonlocal nontrivial interaction between qubits. It seems that difficulties in physical realization of Feynman quantum computers or analogue computers from (Farhi et al. 2000) are significant and some authors even believe (see, for example, (Dyakonov 2003)) that such computers are physically non-realizable. The approach of (Farhi et al. 2000) uses a specific structure of 3-SAT problem. If we use some universal approach that does not take into account a specific form of the problem, we can expect a quantum acceleration in time $`N^{1/2}`$ (Grover 1997), i.e., if a deterministic computer makes a search in $`N`$ steps then quantum universal computers will make the same search in $`N^{1/2}`$ steps. Further, it was shown that acceleration, with a universal approach can be achieved by at most $`N^{1/6}`$ Beals et al. 1998; Preskill 1998. Notice that this quantum $`N^{1/2}`$ acceleration is less than an acceleration that can be obtained by special deterministic and probabilistic algorithms with respect to trivial exhaustive search (Dantsin et al. 2003). Once more quantum machine based on quantum optics phenomena was discussed in (Kazakov 2003). The earlier analogue computational machines were discussed in (Matiyasevich 1987; Blass 1989). In this paper we describe a quantum machine which, as we believe, can be realized physically and which may accelerate solving of some NP-complete problems. Our machine uses different linear and nonlinear quantum optical devices that can be constructed actually and, thus, our computer, at least in principle, can be practically realized. An acceleration that can be obtained heavily depends on different characteristics of our devices. We estimate this acceleration vs. deterministic computers using the characteristics of actually existing devices. ## 2 Statement of problems. Known results. We consider the following problems which are NP-hard (Papadimotriou, Steiglitz 1982; Garey, Johnson 1979). 1 Boolean knapsack, variant 1 Given positive integers $`c_j`$, $`j=1,\mathrm{}n`$ and $`K`$, is there a subset $`S`$ of $`\{0,1,\mathrm{},n\}`$ such that $`_{jS}c_j=K`$? In other words, whether there exist $`n`$ boolean values $`s_i\{0,1\}`$ such that $`_{i=1}^nc_is_i=K\mathrm{?}`$ Here the size $`\|I\|`$ of the instance $`\{n,c_1,c_2,\mathrm{},c_n,K\}`$ is the number of bits needed for binary representations of all the integers $`c_i,K`$. We can suppose, without loss of generality that $`c_i<K`$. Thus, size $`\|I\|`$ can be estimated roughly as $`O(n\mathrm{log}K)`$. 2 Boolean knapsack, variant 2 Given integers $`c_j`$, $`j=1,\mathrm{}n`$ and $`B_{},B_+`$, whether there exist $`n`$ boolean values $`s_i\{0,1\}`$ such that $`_{i=1}^nc_is_i`$ lies in the interval $`(B_{},B_+)`$? Here the instance size is roughly $`O(n\mathrm{log}B_+)`$. 3 Optimization boolean knapsack Given integers $`c_j`$ and $`w_i`$ , $`j=1,\mathrm{}n`$ and the number $`B_+`$, maximize the cost $$\underset{i=1}{\overset{n}{}}w_is_i$$ (1) defined by $`n`$ boolean variables $`s_i\{0,1\}`$ under condition that $$\underset{i=1}{\overset{n}{}}c_is_i<B_+.$$ (2) There is an important difference between the problems 1, 2 and 3. The output in 1,2 is "YES"or "NO". The output of 3 is a number, and one could try to approximate it. Let us remind now some known results about 1, 2 and 3. a All the problems 1-3 can be resolved by exhaustive search in $`2^n`$ steps. b If $`2^n>K`$ then the problem 1 can be resolved by a more effective method, namely, by dynamic programming, in $`O(nK)`$ steps. This method can be described briefly as follows (for details, see Papadimotriou and Steiglitz 1982). The algorithm produces consecutively $`\mathrm{\Sigma }_0,\mathrm{\Sigma }_1,\mathrm{},\mathrm{\Sigma }_n`$ such that $`\mathrm{\Sigma }_j`$ is the set of all possible subsums of $`c_1,\mathrm{},c_j`$ At the first step we set $`\mathrm{\Sigma }_0=\{0\}`$. At $`j+1`$-th step, we set $`\mathrm{\Sigma }_{j+1}=(\mathrm{\Sigma }_j(\mathrm{\Sigma }_j+c_{j+1}))\{0,\mathrm{},K\}`$. The problem has a solution if at the last $`n`$-th step $`K\mathrm{\Sigma }_n`$. In a similar way, we can resolve problem 2, it takes $`O(nB_+)`$ steps, and problem 3, it takes $`O(n\mathrm{min}(B_+,R_{opt}))`$ steps, where $`R_{opt}`$ denotes the maximum cost (1). c Suppose we solve an approximative problem 3, namely, we seek for a number close to a maximal cost. This means that we give up accuracy in exchange for efficiency of our algorithm. To this end we can apply a truncation procedure removing the last $`t`$ digits from the binary representations of $`w_i`$ and $`c_i`$. Let $`w_m`$ be the largest $`w_i`$. Such a procedure leads to a cost $`R_{appr}`$ satifying the estimate (Papadimotriou and Steiglitz 1982; Ibarra and Kim 1975) $$\frac{R_{opt}R_{appr}}{R_{opt}}<ϵ,ϵ=\frac{n2^t}{w_m}$$ (3) The approximating algorithm runs in time $`O(n^4ϵ^1)`$ (Papadimotriou and Steiglitz 1982; Ibarra and Kim 1975). Thus this problem has an approximative solution that can be found in a polynomial number of steps. Notice however, that there are many NP-complete problems that cannot be approximated in such a way, for example, travelling salesman problem (Garey and Johnson 1979; Papadimotriou and Steiglitz 1982). ## 3 Description of the quantum machines Consider $`n+1`$ points $`x_0<x_1<x_2,\mathrm{}<x_n=x_f`$ located along $`x`$ axis in $`(x,y)`$ -plane. At the first point we set a laser, which generates a narrow beam. The diameter of this beam will be denoted by $`d_b`$, the wave length of the laser radiation will be denoted by $`\lambda `$. Typical values of these parameters are $`d_b210^3m,\lambda 510^7m`$. Here we describe an analogue quantum optical device (QOD) for the knapsack problems 1,2. Its possible scheme is presented on fig 1. An input laser beam is splitted by $`50\%`$ mirror A1 to two separate beams running two different trajectories. The beam 1 then is shifted by a plane optical plate on the value $`c_1\kappa `$ in (vertical) z-direction, which is perpendicular to the plane of our picture (here $`\kappa `$ means the minimal shift). Then the beams are on the second mirror B1 and united result (which is a combination of shifted and non-shifted beams) goes through amplifier C1. We suppose that this amplifier has the characteristics presented on fig. 2. After passage of $`m`$ mirrors the propagating light will contain beams shifted in z-direction at all possible distances $`_{i=1}^mc_is_i\kappa `$, where the values $`s_i\{0,1\}`$ depend on the trajectory of the beam. At the final point we measure intensity of outcoming light by a charge-coupled device (CCD). A CCD camera uses a small rectangular piece of silicon to receive incoming light. This silicon wafer is a solid state electronic component which has been micro-manufactured and segmented into an array of individual light-sensitive cells called "pixels". The pixel of the most common CCD has size only about $`\delta _p10^7m0.2\lambda `$, it measures the intensity of light independently from other pixels. Note, that namely $`50\%`$ mirrors are genuine quantum components in our device (Feynman 1985). We denote by $`R_M`$ the z-size of the separating mirrors (and amplifiers). The last parameter important for estimating of quality of our machine is an angular divergence $`\alpha `$ of the laser beam. A usual laser produces a beam having the angular divergence of order $`\alpha d_b/\lambda 410^5`$. Our machine is completely defined thus by the following parameters: $$n,R_M,L,\delta ,Y_{min},Y_{max},\lambda ,d_b,\alpha ,\kappa ,c_i.$$ For the problem 3, we use the following ideas. In order to realize sums (1, 2) we have to modify scheme presented on fig 1. Namely, together with z-shifts we use the horizontal beam shift in an orthogonal direction $`y`$. Notice that a beam shifted in $`y`$-direction contains this shift propagating on different trajectories $`L_1`$ and $`L_2`$ (see fig. 3). It means that after successive z- and y-shifts, which the initial beam gets propagating from $`x_0`$ to $`x_n`$, one obtains the set of beams whose z- and y-shifts are different sums $$\underset{i=1}{\overset{n}{}}c_is_i,\underset{i=1}{\overset{n}{}}w_is_i,$$ in accordance with trajectories. At the end of the device we set CCD which measure the plane distribution of the outcoming light with z-coordinates $`z<B_+\delta _p`$, where $`\delta _p`$ is as it was mentioned above the typical pixel size. ## 4 Solving the dynamic programming by the quantum optical device Consider possible laser beam trajectories as a result of $`n`$ reflections on our mirrors. After the first series of reflections we obtain two beams running along lines $`y=h_1`$ and $`y=0`$ since each photon randomly chooses a way either along the first trajectories or the second one. So, the mirror system serves as a gate choosing the photon way. Remark that the mirror can be considered simultaneously as a quantum and a classical device. In a certain sense, our gates can be thus named as semiclassical quantum gates. If we set a CCD that registrates photons at a point located after $`x_1`$, we could registrate two localized z-separated gaussian beams at $`y=0`$ and $`y=H_1`$. So, in an ideal situation, QOD works as follows. Consider the problem 2. We set $`H_i=\kappa c_i`$, $`i=1,2,\mathrm{},n`$, where $`c_i,n`$ are given. We take now $`Y_{min}=\kappa B_{},Y_{max}=\kappa B_+`$. Now turning on an input laser beam, we check, does the interval $`(Y_{min},Y_{max})`$ contain the center of the any laser beam (respectively, answer either "YES"or " NO"). We suppose that CCD solves the problem with an error less than $`1/3`$. It is not difficult to see that this method of solution can be considered as a physical realization of dynamic programming method from Section 2. The problem 3 can be resolved approximatively, i.e., the global maximum can be estimated within some precision. However, to this end, we must modify slightly our consideration (see below). Consider now diffraction effects. They can create an obstacle since one could registrate a photon not really passing correct gates as a photon that resolves our problem. In accordance with theory of gaussian beams (Svelto 1982) each gaussian beam has a finite transversal size and an angle of divergence $`\alpha `$. These values are connected by $$\alpha \frac{1.2\lambda }{d_b}.$$ (4) After $`n`$ reflections, the size of an output beam becomes $$d_{final}=nL\mathrm{sin}\alpha +d_b,$$ (5) where $`L`$ means distance between mirrors. Therefore, to get minimal value of $`d_{final}`$ we set $$d_b=\sqrt{nL\lambda }.$$ (6) Further, in order to separate on registering CCD the broadening gaussian beams (whose amplitude can differ up to factor 2 in accordance with fig.2), we have to restrict the minimal distance $`\kappa `$ between beams by $$\kappa d_{final}=2d_b.$$ (7) In addition to (4.3) we have a geometric restriction $$R_M>\kappa n+d_b.$$ (8) This inequality means that the shifts of no beams jump out any mirrors. One more problem is a possibility that, at some step $`i=T_0`$ many different beams arrive at the same point. It is not a mathematical difficulty (since this means the existence of many solutions) but it is a physical difficulty since it can lead to a high energy concentration and destruction of the mirrors. To avoid this effect we use active media that saturates the photon density (see fig. 2). This allows us to restrict the photon density at each point of the mirror by some constant. Using of the active media leads however to new difficulty which can give us a lost of solutions. Namely, the phases of the photons of each beam become slightly different after passing the active media. Thus if the problem 2 has many solutions, CCD could not registrate any photons as a result of interference. This effect is possible if a phase shifts are significant. To overcome this difficulty we use special auto phase control devices known in optics. At the points $`z_1,z_2,\mathrm{}z_n`$ (see fig. 1) we places the phase-adjusting devices which measure phase of all beams and compare it with phase of a reference laser beam of the same frequency. Such adjustment gives the possibility to align phases of mixing beams. ## 5 Estimation of QOD performance for boolean knapsack Let us estimate now what we can do using such a QOD and compare this machine performance with deterministic computer performance. ### 5.1 General estimates There are three important parameters: a) the implementation cost; b) the energy cost of solving; c) the running time when we solve $`M`$ times the same problem with different inputs. We compare here the dynamic programming for knapsack problems 1, 2 by a deterministic computer and by QOD. For the sake of unifying the notations we assume, that $`B_+K`$. For the QOD the implementation cost $`CI_{quant}`$ can be estimated roughly as $$CI_{quant}=O(Kn)$$ since the mirror and amplifier sizes are of order $`K`$, the number of the mirrors and the cuvettes is $`n`$. The amplifier size is of order $`Kn`$, one can expect that energetical cost $`CE_{qi}`$ per one instance of the problem also is $$CE_{qi}=O(Kn(n+K)),$$ where factor $`n+K`$ arises due to greatest length of the photon trajectory. The complete energy involved in $`M`$ calculations is then $$CE_{quant}=O(Kn(n+K)M).$$ The implementation time has the order $`Kn`$ and the resolving time $$Time_{quant}leC_1(n+K)M,$$ where $`C_1`$ is inverse proportional to the light speed and thus this coefficient is rather small. For the deterministic computer we have $$CE_{det}=O(KnM),Time_{det}=O(KnM)$$ Thus, QOD wins in time (taking into account both implementation and running times) vs. deterministic machine, while the order of consumed energy is higher than for a deterministic one. ### 5.2 Energetic and temporal costs Our results include energetic and temporal costs of calculations. Let us consider a model, where this correspondence exhibits in a more explicit form. Let several polyhedra be placed in $`𝐑^3`$ and we have to calculate the length of the shortest path from the point A to the point B avoiding these polyhedra (treated as obstacles) (see fig.4). It is well known that this problem is NP-hard (Canny and Reif 1987). Moreover, even determining first $`O(N^{1/2})`$ bits of the length of the shortest path is NP-hard, where $`N`$ denotes the bits size of description of the instance of the problem (Canny and Reif 1987). We can realize this calculation by the following analogue physical device. Let us place in the point A a source of radiation, in the point B - detector and polyhedral obstacles have mirror boundaries. If we turn the source A at initial moment $`t=0`$ and measure the time $`t=t_1`$ when our detector catches a first photon. Let the power of source be $`J`$ and frequency of radiation be $`\nu `$, then the number of created photons is $`J/h\nu `$, where $`h`$ is Planck’s constant (Karlov 1992). Each photon moves in its path and one can describe our analogue calculation as a "parallel computation"of the shortest path by different photons in which each photon plays the role of a processor. In general, this first photon caught by the detector could not belong to the front wave (which corresponds to the shortest paths) due to a weakening of the wave and possible presence of the more strong waves. Therefore, increasing the intensity of the source allows to augment the probability to detect a photon of the wave front. So, we have to either repeat this experiment, with corresponding increasing of temporal cost, or increase the intensity of the source with increasing of energetic cost of calculation. On the other hand, considered analogue devices have brought us to a conjecture that a tradeoff of the following kind $$ETZ,$$ holds, where $`E`$ and $`T`$ are energetical and temporal costs of this calculation respectively, and $`Z`$ depends on the problem. This means (at some extent) that there could be an "exchange"between the energy cost and the temporal cost needed to solve a computational task. ### 5.3 Estimates for real parameters For the rate of a deterministic computer, for a single processor (we do not use a parallelism) we take the value $`V=10^{10}operations/sec`$. Suppose $`2^n>B_+`$, then the dynamic programming is more effective than the exhaustive search. Then a deterministic processor solves the problem 2 within time $`T_c`$ $$T_cnB_+V^1$$ (9) Consider our optical processor. We see that $`B_+`$ must be subject to $$B_+<R_M/\kappa .$$ (10) Thus our possibilities are restricted by the mirror size. Then the problem can be resolved within time $$T_q(nL+R_M)/c+nT_{atom},$$ (11) where the second term describes the time when photons spend in the amplifier (which is at most $`n`$ times the relaxation time $`T_{atom}`$ for active atoms, $`T_{atom}`$ is about $`10^8sec`$). The preprocessing (implementation) time for the optical machine is $$T_{q,pre}C_2Kn,$$ where $`C_2`$ is a constant. It is interesting now to see what we could obtain by using really existing devices. As an example, we take the following typical parameter values: $$R_M=nL=10m,d_b=210^3m,\lambda =510^7m.$$ Then $`\kappa =510^3m`$ and $`B_+<210^2`$. For $`n30`$ the next inequality holds, $`2^n>>nB_+`$, and dynamic programming is more effective than the exhaustive search. Minimal admissible difference $`Y_{max}Y_{min}`$ can be estimated as $`\kappa =510^3m`$. The deterministic time (5.1) is then $`T_{det}30210^210^{10}10^6sec.`$ In accordance with (11) $`T_q30/(310^8)+3010^810^6sec.`$ So, our conclusions for problem 1, 2 are the following. The performance of QOD is restricted by mirror sizes and beam diameter. The problem 2 can be solved by QOD and this effectiveness is essentially more vs. a deterministic processor if we repeat the same computation with different inputs many times. For realistic values of parameters, considered above, we can handle the case $`n<60`$, and then the speeds of the QOD and deterministic processor have the same order. Let us discuss now situation for problem 3. In the next section we describe first modifications needed for resolving this problem. We will see that in this case the speed-up is much better. ## 6 Estimate of machine performance for optimization quantum knapsack ### 6.1 General estimates We compare here the dynamic programming for knapsack problem 3 by a deterministic computer and by QOD using the parameters $`CI`$, $`CE`$, $`Time`$ described in Section 5.1. For the quantum optical device the implementation (preprocessing) cost $`CI_{quant}`$ can be estimated roughly as $$CI_{quant}=O(K^2n),$$ since the the mirror size is now of order $`K^2`$, the number of the mirrors is $`n`$. The amplifier takes volume of order $`K^2n`$, one can expect that the energetical cost $`CE_{quant}`$ is $$CE_{quant}=O(K^2n(n+K)M),$$ where now $`K`$ majorates sums (1) and (2). However, if only numbers $`w_i`$ are large and $`c_i1`$, then the mirror area becomes again $`O(K)`$ as in Section 5, and we have the same estimates as in subsection 5.1. One can estimate the resolving time $$Time_{quant}=O((n+K)M).$$ For the deterministic computer we can (if we solve $`M`$ times the problem with different inputs) estimate energy similarly to 5.1 (problems 1,2) $$CE_{det}=O(KnM),Time_{det}=O(KnM).$$ Thus, if the number of inputs $`M>>K`$ then QOD has a chance to win, in time, with respect to a deterministic processor. The consumed energy of the QOD can be estimated by $$CE_{quant}=O(K^2n(n+K)M)$$ (cf. subsection 5.1 above) whereas the energy consumed by a deterministic machine can be bounded again by $`O(KnM)`$ (cf. the discussion on the tradeoff of energy and time in subsection 5.1). Let us consider now approximating solutions (see Section 2). Recall that the processing time of the deterministic processor in order to obtain an approximative solution, within precision $`ϵ`$, is $$Time_{det,appr}=O(Mn^4ϵ^1).$$ To compare this performance with effectiveness of our QOD, let us note that the pixel size is $`\delta _p`$ and thus the size of the mirrors in QOD solving the problem within relative precision $`ϵ`$ should be $`\delta _p/ϵ=K\kappa `$. Then the implementation cost will be $$CI_{quant}=O((\delta _p/ϵ\kappa )^2n),$$ resolving time can be estimated as $`Time_{quant}=O((n+\delta _p/ϵ\kappa )M)`$ The choice $`\delta _p`$ is restricted by a small number of possible values, thus, $`Time_{quant}+CI_{quant}`$ grows in $`n`$ more slowly than $`Time_{det,appr}`$ for $`ϵ^1=o(n^3M)`$. The consumed energy of QOD for approximating problem can be estimated by $$CE_{quant,appr}=O((\delta _p/ϵ\kappa )^2n(n+\delta _p/ϵ\kappa )M)$$ which is less than the consumed energy for a deterministic machine (which is of the order $`n^4ϵ^1M`$) when $`ϵ^1=o((n\kappa /\delta _p)^{3/2})`$. ### 6.2 Estimates for real device parameters Recall that, to resolve the problem 3, we have modified QOD introducing additional parallel plates performing beam shifts in the direction $`y`$. So, we can suppose that $`i`$-th gate makes the shift to $`H_i`$ in the direction $`z`$ and to $`H_i^{}`$ to the $`y`$ axis. As above, we set $`H_i=\kappa c_i`$ and $`H_i^{}=\kappa w_i`$. At the end of the device we situate a set of CCD, which measure the intensity of light in $`(y,z)`$ plane. Let us describe now how we can resolve, approximatively, any problem 3 without restrictions to the maximum of $`B_+`$. The key truncation idea can be found in (Ibarra and Kim 1975), see also (Papadimotriou and Steiglitz 1982). We reduce the case of arbitrary coefficients $`c_i<B_+,w_i`$ to the previously studied case of restricted coefficients. To proceed it, we use the binary representations of these numbers. Suppose that each number use $`m`$ the digits $`0,1`$, i.e. the size of the binary input is not much than $`m`$. Moreover, assume that the maximal admissible value of $`B_+`$ can be written down with only $`m_{}`$ digits. Now we truncate each given number $`c_i,w_i`$ and $`B_+`$ removing $`mm_{}`$ digits and allowing only first $`m_{}`$ digits. With these new truncated data we can solve our problem and this solution gives an approximative solution. The solving procedure to search an approximative solution of the problem 3 can be described as follows. We observe all beams that are measured by CCD and between them we choose a beam that have a maximal y-deviation. Such a scheme always gives an approximative solution of our problem, with precision of order $`\kappa /2`$. Let us compare now a deterministic processor and QOD. Let us take the same parameters as above. The time processing for our device will be chosen the same as above, i.e., $`T_q=10^6sec`$. For the deterministic computer, according to subsection 2c, one has $$T_cn^4Vϵ^1,$$ where $`ϵ`$ is the relative precision. The relative error of measurement we can estimate as $`ϵ(\kappa /2)/R_M`$, which corresponds to the separation of the gaussian beams with amplitudes differs by factor like 2. Substituting all the values in the last relation, we have $`T_c1sec`$, that much more than $`T_q`$. So, the quantum optical device gains an acceleration $`10^6`$ times vs a typical deterministic computer. ## 7 Conclusion We have introduced and designed analogous quantum optical devices for simulating dynamic programming which provides the following complexity bounds (see notations in subsections 5.1, 6.1) for different versions of the knapsack problem (section 2). ###### Proposition 7.1. i) For the versions 1,2 the implementation cost $`CI_{quant}=O(Kn)`$, the running time $`Time_{quant}=O((n+K)M)`$, and the energy cost $`CE_{quant}=O(Kn(n+K)M)`$; ii) For the version 3 $`CI_{quant}=O(K^2n)`$, the running time $`Time_{quant}=O((n+K)M)`$, and the energy cost $`CE_{quant}=O(K^2n(n+K)M)`$; iii) for $`ϵ`$-approximative solution of the version 3 $`CI_{quant}=O(n(\delta _p/ϵ)^2)`$, the running time $`Time_{quant}=O((n+\delta _p/ϵ)M)`$, and the energy cost $`CE_{quant}=O(n(n+\delta /ϵ)(\delta _p/ϵ)^2M)`$; (see subsections 5.1 and 6.1) Also we compare these bounds with ones for deterministic machines (see subsections 5.1, 6.1). Let us discuss briefly the "quantum properties"of QOD. The genuine quantum machine has to exploit two quantum properties: i) "exponential resource"connected with exponentially large dimension of the state space, and ii ) the "quantum parallelism"which is simultaneous evolution in all subspaces of the state space. The QOD described above contains only one genuine quantum component, namely, $`50\%`$ mirrors, which split the laser beams. This splitting gives the possibility to realize the "quantum parallelism". The splitting mirrors operate with the laser beams that are macroscopic objects. However, this macroscopity prevents to the realization of "exponential resource". In this framework, realization of the exponential resource means using of the mirrors of exponential size in order to separate the exponentially large number of the laser beams. Moreover, in this case we need exponentially large energetic resource in order to keep the intensity of the laser beams on a suitable level. So, it is difficult to realize of the "exponential resource with help of macroscopic quantum devices. But the second quantum property - "quantum parallelism can be realized by described above devices. ## 8 Ackowledgements The second and third authors are grateful to the Mathematical Institute of the University of Rennes, France, for the hospitality.
warning/0506/cond-mat0506204.html	ar5iv	text	# Generalized coarse graining applied to a ϕ⁴-theory: A model reduction-renormalization group synthesis ## Abstract We develop an algorithmic, system-specific renormalization group (RG) procedure by implementing a systematic coarse graining procedure that is adapted from model reduction techniques from engineering control theory. The resulting “generalized” RG is a consistent generalization of the Wilsonian RG. We apply the generalized RG to a $`\varphi ^4`$ field theory. By considering nonequilibrium in addition to equilibrium observables, we find that naïve power counting breaks down. Additionally, a large class of short-wavelength perturbations can drive the system away from both the Gaussian and Wilson-Fisher fixed points. Although it is already known that the renormalization group (RG) is not a black box routine, the purpose of this letter is to make it more algorithmic. It is easy to be misled into thinking that the RG already is algorithmic because its key ingredients are coarse graining and rescaling the system variables goldenfeld92 ; shankar94 . However, fully algorithmic implementations of the RG fail for large classes of problems because it is not possible to ignore the physics of a system and expect to obtain meaningful results. Capturing the essential physics requires isolating the appropriate models and the structure of perturbations and uncertainties. Consequently, this process is system-specific. Additionally, the scale on which the physics is observed must be specified. For instance, for bosonic theories the long-wavelength physics is investigated, while for fermionic system it is the physics near the Fermi surface. These considerations suggest that the primary obstacles to automation are the model identification and coarse graining processes. In this article we present a RG procedure that accounts for these system-specific obstacles and apply it to investigate the equilibrium properties of a $`\varphi ^4`$-theory. The perspective adopted in this letter is that the obstacles mentioned above arise from the identification of physical observables. We use control theoretic techniques for open systems to systematically identify observables. Specifically, consider an open system of the form $$\dot{𝐱}=𝐟(𝐱)+𝐮,$$ (1) where $`𝐱`$ and $`𝐮`$ are vectors in possibly infinite dimensional spaces. From a field theoretic or statistical mechanical point of view, without $`𝐮`$, equation (1) represents the “classical” equations of motion of the system. $`𝐮`$ represents generic driving as well as a possible noise to which the system is exposed. If $`𝐮`$ is generated by a continuous stochastic process, this imposes a specific structure on the noise. Theories with this restriction on $`𝐮`$ may be mapped to a field theory via Martin-Siggia-Rose (MSR)/closed-time-path (CTP) methods martin73 ; cooper01 . In these instances, coarse graining is typically dictated by conserved quantities and thermodynamic considerations. Consequently such systems are locally coarse grained. Suppose $`𝐮`$ is an arbitrary input (e.g. $`𝐮L_2`$) to the system. We consider the states or regions in phase space that are most accessible via driving to be responsible for describing the the essential characteristics of the system. This is analogous to the energy landscape picture in statistical mechanics where fluctuations govern which states contribute the most to the statistics of the system. We use this control theoretic notion of the importance of states to specify how to coarse grain and, consequently, to generate RG equations. The key step in generalizing the RG lies in ascertaining how to coarse grain. While equation (1) addresses the effects of perturbations to the system, it does not allow for the possibility of multiscale or constrained observation. We can remedy this by considering more general open systems of the form $$\begin{array}{c}\dot{𝐱}=𝐟(𝐱)+B𝐮,\\ 𝐲=C𝐱,\end{array}$$ (2) where $`𝐲`$ reflects that only some subspace or, more generally, subset of phase space is directly measurable. The operators $`B`$ and $`C`$ respectively specify the structure of how noise may enter the system and which states can be measured. The additional structure provides a practical way to model real systems and consider experimental constraints. The constraint imposed by only measuring $`𝐲`$ strongly influences the relative importance of the internal states $`𝐱`$ and hence coarse graining. For instance, if we measure a projected subspace of $`𝐱`$ over a finite but short time horizon, we can expect that those states responsible for the transient dynamics will be the most important. In this case, conservation laws may play an insignificant role in determining how to coarse grain. In reynolds04 we proposed to coarse grain linear systems based on retaining the states contributing most to the response of the system to disturbances. Following common practice, we coarse grain equation (2) based on its linearization about a particular solution with $`𝐮=0`$. In particular, to simplify analysis, we only consider linearizations about equilibrium solutions. The linearizations are generically described by $$\begin{array}{c}\dot{\stackrel{~}{𝐱}}=A\stackrel{~}{𝐱}+\stackrel{~}{B}\stackrel{~}{𝐮},\\ \stackrel{~}{𝐲}=C\stackrel{~}{𝐱}.\end{array}$$ (3) Associated with equation (3) are invariants known as Hankel singular values (HSV’s) glover84 ; peller03 . The HSV are nonnegative real numbers, $`\sigma _{\mathrm{max}}\sigma _\kappa \sigma _{\mathrm{min}}`$ that comprise the spectrum of the operator $`W`$. For systems considered over an finite time horizon, $`t_f`$, consider the positive operators $`X`$ and $`Y`$, called gramians, that are determined by the equations $`{\displaystyle \frac{\mathrm{d}X}{\mathrm{d}t_f}}=AX+XA^{}+\stackrel{~}{B}\stackrel{~}{B}^{};X(0)=0,`$ (4) $`{\displaystyle \frac{\mathrm{d}Y}{\mathrm{d}t_f}}=A^{}Y+YA+C^{}C;Y(0)=0.`$ (5) $`W^2`$ may be factored as $$W^2=XY.$$ (6) HSV’s provide a precise measure of the error incurred by approximating the effect $`\stackrel{~}{𝐮}`$ has on $`\stackrel{~}{𝐲}`$ with reduced order models. The HSV’s may be interpreted as supplying a measure of the importance of the internal states $`\stackrel{~}{𝐱}`$. If $`W`$ is invertible, it is always possible to find a coordinate system, called balanced coordinates, such that $`X=Y=\mathrm{diag}(\sigma _{\mathrm{max}},\mathrm{},\sigma _{\mathrm{min}})`$ <sup>1</sup><sup>1</sup>1Some subtleties arise for infinite dimension systems, but otherwise similar results hold.. When equation (3) is transformed to balanced coordinates, the best reductions are those that project out the states corresponding to small HSV. In other words, the ordering of the HSV, at least locally around an equilibrium configuration in phase space, specifies how to coarse grain a system. An in depth treatment of this material may be found in reynolds04 ; dullerud00 . It is also sometimes possible to “balance” the full nonlinear system scherpen93 . The RG can easily be adapted for HSV-based coarse graining. Operator theoretic approaches to the RG muller96 ; bach98 demonstrate that coarse graining in the Wilsonian RG is equivalent to multiplying operators or states by projection operators <sup>2</sup><sup>2</sup>2As is indicated in bach03 , we need not limit our attention to projection operators.. The essense of this work is to use HSV’s to identify the projection operator. As before, suppose that $`\kappa `$ is a vector index that orders the HSV’s $`\sigma _\kappa `$ for equation (3) from largest to smallest. A generalized Wilsonian RG procedure is obtained by: 1) transforming the system to balanced coordinates (about an equilibrium solution), $$\widehat{\varphi }(\kappa ,t)=R(\kappa ,𝐱)\varphi (𝐱,t)d𝐱,$$ (7) so that the partition function takes the form, $$\begin{array}{c}𝒵\left[\{\varphi \}\right]=𝒟\varphi \mathrm{exp}\left(S(𝐠,\{\varphi \})\right)\hfill \\ =𝒟\widehat{\varphi }𝒥\mathrm{exp}\left(S(\stackrel{~}{𝐠},\{\widehat{\varphi }\})\right),\hfill \end{array}$$ (8) where $`𝐠`$ is the original set of coupling constants/functions, $`𝒥`$ is the Jacobian from equation (7), and $`\stackrel{~}{𝐠}`$ is the resulting transformed set of coupling constants; 2) integrating out $`\kappa `$-shells about $`\sigma _{\mathrm{min}}`$ analogously to wavevector shells; and 3) rescaling $`\kappa `$ and $`\widehat{\varphi }`$ appropriately. An interesting but technically challenging variant of this procedure is to integrate out $`\sigma _\kappa `$-shells instead of $`\kappa `$-shells about $`\sigma _{\mathrm{min}}`$. The remarkable feature of this variant is that $`\sigma _\kappa `$ does not respect spatial dimension. For instance, integrating out a single $`\sigma _\kappa `$-shell may entail integrating out an entire subspace in $`𝐱`$-space. This seems to provide a natural way to understand how three dimensional systems may have regions (e.g. affine subspaces) exhibiting one or two dimensional critical behavior. The technical challenge lies in rescaling $`\sigma _\kappa `$. It is not clear that rescaling $`\sigma _\kappa `$ will recover the full $`\kappa `$-space thereby generating a meaningful RG. Before applying this procedure to the nonlinear wave equation, we apply it to some trivial examples to build intuition. We first consider the (driven) diffusion equation $$_t\varphi =D^2\varphi +\gamma u.$$ (9) In this example, $`B=\gamma `$, $`C=1`$, and we take $`t_f\mathrm{}`$. By considering a stable system over an infinite time horizon, we only need to solve the Lyapunov equations, $`AX+XA^{}+\stackrel{~}{B}\stackrel{~}{B}^{}=0,`$ (10) $`A^{}Y+YA+C^{}C=0,`$ (11) instead of equations (4)-(5). By taking the Fourier transform of equations (10)-(11) it is easy to derive that $`W`$, from equation (6), in balanced coordinates is given by $$W_𝐤^{bal}=\frac{\|\gamma \|}{2D\|𝐤\|^2}.$$ (12) Here the index, $`\kappa `$ for the HSV’s is just $`\|𝐤\|`$. Thus, for the diffusion equation, the most important states are those that correspond to small wavevector. Thus, local coarse graining is appropriate because the smallest observable “fluctuations” are due to the short-wavelength physics. The smallest error is incurred by projecting out large wavevectors. We will not complete the RG analysis here because the standard RG treatments based on local coarse graining are applicable. For a field theoretic treatment it is possible to obtain the coarse grained action from the MSR/CTP formalism zanella02 , however alternative formulations may be found in goldenfeld92 ; bricmont95 . As a second example, we consider the (driven) linear wave equation $$\begin{array}{c}_t^2\varphi =v^2^2\varphi +\gamma u,\\ 𝐲=\varphi .\end{array}$$ (13) By the units of $`u`$, it represents a true force acting on $`\varphi `$. This, in addition to the fact that only $`\varphi `$ is the “measurable” quantity, implies that we have isolated our attention on $`\varphi `$-based observables. This is choice is a very statistical equilibrium and thermodynamic one. We have completely neglected $`\pi `$, the field conjugate to $`\varphi `$, that represents the kinetic contributions to the system. When posed as a set of first order equations, equation (13) becomes $$\left[\begin{array}{c}_t\varphi \\ _t\pi \\ 𝐲\end{array}\right]=\left[\begin{array}{ccc}0& 1& 0\\ v^2^2& 0& \gamma \\ 1& 0& 0\end{array}\right]\left[\begin{array}{c}\varphi \\ \pi \\ u\end{array}\right]$$ (14) By smoothing out the time-cutoff at $`t_f`$ with a damped exponential in the integral representation of the solution of equations (4) and (5) the problem simplifies to solving Lyapunov equations. This smoothing process is also known as exponential discounting. With the given form of $`B`$ and $`C`$ in this problem, we find that the matrix of HSV’s, $`W`$, is approximately given by $$W_𝐤^{bal}\frac{\|\gamma \|}{4av\|𝐤\|}I_{2\times 2},$$ (15) where $`I_{2\times 2}`$ is the $`2\times 2`$ matrix identity, $`a1/t_f`$, and $``$ is the dyadic (algebraic tensor) product. As with the diffusion equation, short-wavelength physics does not significantly contribute to the response, so locally coarse graining is appropriate. we will refrain from treating this example in more detail because the standard RG treatment of the wave equation and its $`\varphi ^4`$ nonlinear generalization may be readily found in field theory textbooks peskin95 . The main point of the previous examples is to convey that the generalized Wilsonian RG consistently reproduces the results of the standard RG for systems that we already intuitively know should be locally coarse grained. A great strength of the generalized RG is that it permits us to tackle less intuitive problems. We apply the generalized RG to determine the statistical equilibrium properties of a nonlinear wave equation with a particular nonequilibrium choice of observables. As will be seen, a surprising result is that this choice of observables forces us to nonlocally coarse grain. The nonlocality of the coarse graining has very interesting implications with regard to the resulting induced RG flow. In the remainder of the paper, we then consider some of these novel implications. The (driven) equations of motion that we are considering are $$\begin{array}{c}_t\varphi =\pi +\alpha _1u_1\hfill \\ _t\pi =^2\varphi +\frac{\lambda }{3!}\varphi ^3+\alpha _2u_2\hfill \\ 𝐲=\left[\begin{array}{c}\beta _1\varphi \\ \beta _2\pi \end{array}\right],\hfill \end{array}$$ (16) where $`\varphi `$ and $`\pi `$ are real-valued fields. The driving now includes generalized forces in addition to “true” forces. By expanding around equilibrium solutions of $`^2\varphi =0`$ we find that for each real-space position $`𝐱`$, $`B=\left(\begin{array}{cc}\alpha _1& 0\\ 0& \alpha _2\end{array}\right)\text{and}`$ (19) $`C=\left(\begin{array}{cc}\beta _1& 0\\ 0& \beta _2\end{array}\right).`$ (22) This driving allows for more states in $`(\varphi ,\pi )`$-phase space to be accessible compared to the driving in equation (14). This, in combination with the form of $`𝐲`$, ensures that both $`\varphi `$ and $`\pi `$-dependent observables are being considered. By using exponential discounting, as mentioned earlier, we find that the diagonal operator of HSV’s is given by $$\begin{array}{c}W_𝐤\frac{1}{4a}[(\alpha _2^2\|𝐤\|^1+\alpha _1^2\|𝐤\|)\\ \times (\beta _1^2\|𝐤\|^1+\beta _2^2\|𝐤\|)]^{1/2}I_{2\times 2}.\end{array}$$ (23) $`W_𝐤`$ does not have the HSV’s ordered from largest to smallest, so it is not expressed in balanced coordinates. It is immediately apparent that the HSV’s are large for both large and small magnitude wavevector. A heuristic explanation for this strange result is that for large wavevector, $`\pi `$ is a pathologically “fast” variable. However, by driving $`\pi `$ with $`u_1`$ over all wavevector, this permits the fast resonances to be excited at large wavevector. The pathological nature of $`\pi `$ as an observable is analogous to the pathological nature of considering $`\dot{\xi }`$ an observable where $`\xi `$ satisfies a Langevin equation <sup>3</sup><sup>3</sup>3$`\dot{\xi }`$ is pathologically fast compared to $`\xi `$. For this reason, $`\pi `$ is a nonequilibrium observable. Furthermore, because both the small and large wavelength physics contributes strongly to the response of the system, local coarse graining cannot be correct. The appropriate coarse graining is nonlocal. In the case where $`\varphi `$ and $`\pi `$ are treated on equal footing as observables, which in general may not be the case, $`\alpha _1=\alpha _2=\alpha `$ and $`\beta _1=\beta _2=\beta `$. In the remainder, we treat this particular case. Furthermore, without loss of generality, we set $`\alpha =\beta =1`$. In this case, equation (23) indicates that the $`\|𝐤\|=1`$ states are the least important. Implementing the second step of the procedure for generalized RG involves integrating out $`𝐤`$-shells away from the $`\|𝐤\|=1`$ surface. Rather than transform the system into the balanced $`\kappa `$-coordinates, out of convenience, we coarse grain the system in $`𝐤`$-space. The action for this system, without driving, in Fourier space is given by $$\begin{array}{c}S(𝐠,\{\widehat{\varphi }\})=\frac{1}{(2\pi )^D}d𝐤\|𝐤\|^2\widehat{\varphi (𝐤)}^2+\hfill \\ \frac{\lambda }{4!}_{n=1}^4\frac{\mathrm{d}𝐤_n}{(2\pi )^D}\delta \left(_{j=1}^4𝐤_j\right)\widehat{\varphi }(𝐤_1)\widehat{\varphi }(𝐤_2)\widehat{\varphi }(𝐤_3)\widehat{\varphi }(𝐤_4),\hfill \end{array}$$ (24) where $`D`$ is the spatial dimension of the system since we are only considering the statistical equilibrium properties of the system. If we wished to do a full nonequilibrium treatment, we should coarse grain the MSR/CTP-action for the system. In order to coarse grain, we let $`\widehat{\varphi }=\widehat{\varphi }_<+\widehat{\varphi }_m+\widehat{\varphi }_>`$ where $`\widehat{\varphi }_<`$ is only nonzero for $`\|𝐤\|\mathrm{\Lambda }`$, $`\widehat{\varphi }_m`$ is only nonzero for $`\mathrm{\Lambda }<\|𝐤\|<\mathrm{\Lambda }^1`$, and $`\widehat{\varphi }_>`$ is only nonzero for $`\|𝐤\|\mathrm{\Lambda }^1`$, where $`\mathrm{\Lambda }<1`$. With this decomposition, the path integral measure factors as $`𝒟\widehat{\varphi }=𝒟\widehat{\varphi }_<𝒟\widehat{\varphi }_m𝒟\widehat{\varphi }_>`$. The RG equations are then induced by integrating out $`\widehat{\varphi }_m`$ and then rescaling the wavevectors and fields. For this problem, the rescaling procedure requires that $$\begin{array}{c}\widehat{\varphi }_<(𝐤)=Z_<\phi _1(\mathrm{\Lambda }^1𝐤),\\ \widehat{\varphi }_>(𝐤)=Z_>\phi _2(\mathrm{\Lambda }𝐤),\end{array}$$ (25) and $`𝐩=\mathrm{\Lambda }^1𝐤`$ for $`\|𝐤\|\mathrm{\Lambda }`$ and $`𝐩=\mathrm{\Lambda }𝐤`$ for $`\|𝐤\|\mathrm{\Lambda }^1`$. Naïve power counting breaks down as a direct result of rescaling in the two disjoint wavevector regimes. Although we start with a theory where $`𝐠=(1,\lambda ,0,\mathrm{})`$ we can expect that the RG transformations may generate new nonlinear terms and that the coupling constants may become coupling functions. In fact, $`𝐠`$ flows towards having an infinite number of nontrivial components. In particular, the coupling constant $`\lambda `$ becomes a coupling function, $`\lambda (𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐},𝐩_\mathrm{𝟑},𝐩_\mathrm{𝟒})`$, that may be decomposed into the coupling functions $`\{\lambda _{(i,4i)}\}_{i=1}^4`$. Here the notation indicates that $`\lambda _{(i,j)}`$ has $`i`$ wavevectors with $`\|𝐩_n\|<1`$ and $`j`$ wavevectors with $`\|𝐩_n\|>1`$. If we let $`\mathrm{\Lambda }=e^{dl}`$, then to first loop order the RG equations for $`\lambda _{(i,j)}(𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐},𝐩_\mathrm{𝟑},𝐩_\mathrm{𝟒})`$ are $$\begin{array}{c}_l\lambda _{(i,j)}=\alpha _{i,j}\lambda _{(i,j)}\hfill \\ +4^1\lambda _{(i,j)}^2J(𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐},𝐩_\mathrm{𝟑},𝐩_\mathrm{𝟒})+𝒪(\lambda ^3),\hfill \end{array}$$ (26) where $`\alpha _{0,4}=D4`$, $`\alpha _{4,0}=4D`$, $`\alpha _{2,2}=D`$, $`\alpha _{3,1}=2(D1)`$, and $`\alpha _{1,3}=2`$. Also, $`J(𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐},𝐩_\mathrm{𝟑},𝐩_\mathrm{𝟒})=`$ $`K(𝐩_\mathrm{𝟏}+𝐩_\mathrm{𝟐})+K(𝐩_\mathrm{𝟏}+𝐩_\mathrm{𝟑})+\mathrm{}+K(𝐩_\mathrm{𝟑}+𝐩_\mathrm{𝟒})`$, where $`K(𝐪)`$ comes from the one-loop contribution to the 4-point vertex and is given by $$\begin{array}{c}K(𝐪)=\frac{2}{\mathrm{\Lambda }\mathrm{\Lambda }^1}_{\mathrm{\Lambda }<\|𝐤\|<\mathrm{\Lambda }^1}\frac{\mathrm{d}𝐐}{(2\pi )^D}\frac{1}{𝐐^2(𝐐+𝐪)^2}\\ =\frac{2^{D1}S_{D1}}{(2\pi )^D\left(1+\|𝐪\|\right)^2}B(\frac{D1}{2},\frac{D1}{2})\\ \times {}_{2}{}^{}F_{1}^{}(D1,\frac{D1}{2};1;\frac{4\|𝐪\|}{(1+\|𝐪\|)^2}),\end{array}$$ (27) where $`B(x,y)`$ is the Euler beta function and $`{}_{2}{}^{}F_{1}^{}`$ is the hypergeometric function. Although we do not start with a “mass” term in the action (i.e. $`m^2\varphi ^2`$), such a term is generated by the RG flow. If we denote the mass terms for $`\|𝐪\|<1`$ and $`\|𝐪\|>1`$ respectively by $`m_<^2(𝐪)`$ and $`m_>^2(𝐪)`$, then the associated RG equations for these terms are given by $`_lm_<^2=2m_<^2+{\displaystyle \frac{S_D\lambda }{2(2\pi )^D}}+𝒪(\lambda ^2)`$ (28) $`_lm_>^2=2m_>^2+{\displaystyle \frac{S_D\lambda }{2(2\pi )^D}}+𝒪(\lambda ^2)`$ (29) The first thing to notice in equation (26) is that the contribution from tree level, the linear term, indicates that the coupling functions involving a mixing of wavevectors (i.e. $`i,j0`$) are irrelevant. This is the first piece of evidence that the theory decouples at small and large wavelength. The second piece of evidence for this is that $`K(𝐪)`$ in equation (27) diverges as $`\|𝐪\|1`$. Due to $`K(𝐪)`$-dependence $`\lambda `$ acquires through $`J(𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐},𝐩_\mathrm{𝟑},𝐩_\mathrm{𝟒})`$, this divergence seems to act like a barrier in the effective coarse grained theory to prevent the large and small wavevector physics from mixing. As alluded to earlier, this decoupling of the small and large wavelength physics is a manifestation of independence of $`\varphi `$ and $`\pi `$ as physical observables. Rather than being general, we consider this decoupling to be special because we are only considering a real, scalar $`\varphi ^4`$-theory. Without driving, this model lacks continuous symmetries to be broken. In a more general model, the existence of such symmetries and the associated set of gauge transformations would provide means of coupling $`\varphi `$ and $`\pi `$. It is not possible to ignore the wavevector dependence that $`\lambda `$ acquires because they are relevant for $`\|𝐩_i\|>1`$. Specifically, higher derivative perturbations, $`p^n\widehat{\varphi }^2,n>2`$ and $`p^n\widehat{\varphi }^m,n>0,m4`$, in addition to higher order nonlinearities, $`\widehat{\varphi }^n,n>4`$, become relevant when $`\|𝐩_i\|>1`$. However, the couplings at small wavevector, $`\|𝐩_i\|>1`$, obey the standard RG equations obtained by local coarse graining. Physically this means that just as long as the system is not exposed to short-wavelength perturbations, the long-wavelength, “thermodynamic” physics will remain robustly observable. If the system is perturbed to its large-wavevector regime, then it will flow to a short-wavelength fixed point instead of the more familiar Gaussian and Wilson-Fisher fixed points. This reflects that the dynamically faster short-wavelength perturbations are able to excite the conjugate field $`\pi `$, thereby driving the system away from its standard statistical equilibrium. Were the conjugate field not accessible to the “noise”, $`\alpha _1=0`$, or not an observable, $`\beta _2=0`$, this phenomena would not have occurred. In this letter we have presented a new RG procedure and have applied it to a $`\varphi ^4`$ toy model. This procedure provides, to the author’s knowledge, the first systematic means to identify the RG projection operator. When both equilibrium and nonequilibrium observables are chosen, this RG procedure predicts that naïve power counting breaks down and that terms that are ordinarily irrelevant become relevant at large wavevector. Table 1 summarizes these results for the $`\varphi ^4`$-theory. The generalized Wilsonian RG developed here is applicable to nonequilibrium and heterogeneous systems, finite or infinite dimensional systems, and systems with various perturbations and uncertainties. Although the RG is still formally an uncontrolled approximation, the coarse graining is chosen such that the effective, coarsened system is close to the original one. Despite the versatility of this method, it is often difficult to analytically determine the balancing transformations. However, since there are very efficient numerical algorithms for finding balanced coordinates, this generalized RG remains a numerically useful and practical algorithm. ###### Acknowledgements. This work was supported by NSF Grant No. DMR-9813752. Special thanks are due to Jean Carlson for her valuable comments.
warning/0506/astro-ph0506434.html	ar5iv	text	# Formation and Evolution of Planetary Systems: Cold Outer Disks Associated with Sun-like stars ## 1 INTRODUCTION Studying the formation of our own solar system and observing the frequency of similar systems associated with other stars are two ways in which we seek to understand our origins. Through remote observation and direct exploration, we have developed a much clearer understanding of our solar system. However many questions about the processes involved in the initial formation and subsequent evolution towards the present configuration cannot be addressed directly. Therefore, we need to study other stars to help place our solar system in context. There are two major zones of debris in the solar system: the asteroid belt at 2 $``$ 4 AU composed of rocky material that is ground up by collisions to produce most of the zodiacal dust cloud and the Kuiper Belt (KB) that consists of small bodies orbiting beyond Neptune’s orbit at 30 $``$ 50 AU. Since their discovery over a decade ago (Jewitt & Luu 1993), Kuiper Belt objects have played an increasingly important role in understanding the formation and evolution of our planetary system (e.g. Malhotra 1993; Kenyon & Bromley 2004). While direct detection of in situ debris from collisions amongst Kuiper Belt objects has yet to be confirmed, of order 10% of the solar system’s far-IR luminosity could be emitted by Kuiper Belt dust (e.g. Backman et al. 1995). Searching for KB-like debris around other Sun-like stars as a function of age will determine both the frequency of such systems and provide important insight into the formation of our solar system. Most debris disks were found by their mid-to-far infrared emission in excess of the expected photosphere (e.g., Auman et al. 1984; Backman & Paresce 1993; Mannings & Barlow 1998; Habing et al. 1999; Silverstone 2000; Spangler et al. 2001, Decin et al. 2003, Zuckerman & Song 2004). The majority of these debris disks are also associated with hot, luminous stars, since observatories such as IRAS and ISO did not have the necessary sensitivities to detect debris disks around lower-luminosity solar-type stars at distances beyond a few parsecs. The increased sensitivity afforded by the Spitzer Space Telescope (Spitzer; Werner et al. 2004) has the potential to identify and investigate debris systems that were not detectable with previous observatories. Using data from the Formation and Evolution of Planetary Systems (FEPS) Spitzer Legacy Program, we have conducted a search for outer dust disks dominated by temperatures characteristic of the Kuiper Belt ($`T3060`$ K). The earliest results from our validation observations are presented in Meyer et al. (2004), where we identified such debris disks (exo-KBs) surrounding the 30 Myr-old Sun-like star HD 105 and the 0.4 - 1 Gyr old star HD 150706. Herein we present five more Sun-like stars that exhibit characteristics of exo-KBs, and discuss their properties in the context of the evolution of our own solar system. We also present one star that has a firm detection at 70 $`\mu `$m, but the measurement is consistent with the photospheric emission from the star. In this case, we place unprecedentedly low upper limits on the presence of KB-like debris in this system. In § 2, we describe the observations and data reduction. In § 3, we discuss the methodology for identifying the five exo-Kuiper Belt candidate stars from the FEPS sample, and we present the resulting spectral energy distributions (SEDs). We outline our interpretation of the SEDs in terms of physical models in § 4. We discuss the implications of our results in § 5, and summarize our findings in § 6. ## 2 OBSERVATIONS & DATA REDUCTION The FEPS program is described by Meyer et al. (2004, 2005) and a detailed description of the data acquisition and data reduction for the FEPS program is given in the FEPS V1.1 Explanatory Supplement (Hines et al. 2004). In this section, we present a condensed description as it applies to the five targets investigated herein. Observational details are listed for each star in Table 1. Imaging data were obtained with the Multi-band Imaging Photometer for Spitzer (MIPS: Rieke et al. 2004) in the 24 $`\mu `$m and 70 $`\mu `$m bands using small field photometry mode with 2 – 10 cycles of 3 and 10 second integration times respectively; each of such dataset is a data collection event (DCEs). After initial processing by the Spitzer Science Center (SSC) version S10.5.0 pipeline to provide reconstructed pointing information, we used the MIPS Data Analysis Tool (DAT) software (Gordon et al., 2004, 2005) to process the data. For the 24 $`\mu `$m data, read-2 correction, dark subtraction, droop correction, electronic nonlinearity correction, scan mirror dependent flat fields, cosmic ray rejection, and distortion correction were applied. We produced both individual DCE frames and a mosaic image using all DCE frames. Outlier rejection of each pixel was performed for a mosaic image. The typical number of total DCE frames per source used for photometry was 28 for 2 cycles. For the 70 $`\mu `$m data, dark subtraction, illumination correction, and cosmic ray identification/rejection were applied. The electronic nonlinearity of the detector was corrected. (Gordon et al., 2004, 2005) note in their investigation that this nonlinearity is $``$ 1 %, much smaller than absolute calibration uncertainties of 20 % (see the Spitzer Observers Manual 4.6: SOM 4.6, and the MIPS Data Handbook 2004). A “time-filtering” algorithm as described by (Gordon et al., 2004, 2005) was used to eliminate time-dependent gain drifts. All 70 $`\mu `$m flux measurements were performed on the distortion corrected final mosaic image. We used the IDL-based software package IDP3 (Schneider & Stobie 2002) to perform aperture photometry on the 24$`\mu `$m and 70$`\mu `$m data. IDP3 was developed by the Instrument Definition Team for the Near Infrared Camera and Multi-object Spectrometer (NICMOS), and has been further optimized for the Spitzer IRAC and MIPS data. Object fluxes were measured using standard aperture photometry techniques. For the 24 $`\mu `$m photometry, we adopted a 14$`\stackrel{}{\mathrm{.}}`$7 aperture radius and a background annulus from 29$`\stackrel{}{\mathrm{.}}`$4 to 41$`\stackrel{}{\mathrm{.}}`$7 to perform aperture photometry on the distortion corrected individual DCE frames. For the 70 $`\mu `$m photometry we used a target aperture radius of 29$`\stackrel{}{\mathrm{.}}`$7, and a background annulus of 39$`\stackrel{}{\mathrm{.}}`$6 – 79$`\stackrel{}{\mathrm{.}}`$2 on the single mosaic image. Aperture correction factors, listed in the SOM 4.6, were then applied. For HD 8907, we used a smaller target aperture (19$`\stackrel{}{\mathrm{.}}`$7) with the same background annulus at 70$`\mu `$m to avoid contamination by a nearby source. We also placed a mask on the nearby source to minimize any contribution from the wings of the PSF to the background annulus. The aperture correction for this smaller aperture was derived relative to the standard aperture using 70um observations from other bright and isolated FEPS stars. The uncertainty in this small-aperture correction is $``$ 1.2%. The 24$`\mu `$m flux uncertainties are reported in Table 2, and include both the standard deviation of the mean flux density for individual DCEs, which is an estimate of the internal precision of our measurement, and the absolute calibration uncertainty. The 70$`\mu `$m internal uncertainty is estimated by the rms pixel-to-pixel dispersion inside the background annulus measured on the mosaic image and scaled to the area of the target aperture. The total uncertainty for both bands is the combination of the internal uncertainty combined (i.e., added in quadrature) with the calibration uncertainty. The calibration uncertainties for 24 $`\mu `$m and 70 $`\mu `$m photometry are 10 % and 20 % respectively (see the Spitzer Observers Manual 4.6 and MIPS Data Handbook 2004). Note that the the 70 $`\mu `$m internal uncertainties are much smaller than absolute calibration uncertainty, and the internal signal-to-noise is very high for our detected objects. It is our knowledge of the absolute flux density in physical units that retains the large uncertainties in Table 2. Infrared Array Camera (IRAC: Fazio et al. 2004) observations at 3.6, 4.5 and 8.0$`\mu `$m were also obtained for each object using the 32$`\times `$32 pixel subarray mode with an effective integration time of 0.01 sec per image (frame time of 0.02 s). Sixty four images of the object were obtained at four different positions in a random dither pattern on the array. This gives a total of 256 images (2.56 sec total integration time) per band. We used the basic calibrated data (BCD) products produced by the SSC S10.5.0 data pipeline as described in the SOM 4.6. Aperture photometry was performed with IDP3 using a 3 pixel radius aperture centered on the target. The background was estimated using the 820 pixels that lie beyond an 10 pixel radius relative to the star (i.e., all pixels beyond the target radius). The background flux was normalized to the area of the target aperture and subtracted from the summed target flux. The final source flux is the mean of the 256 measures, corrected from a 3 pixel radius to the 10 pixel radius used for the IRAC instrumental absolute flux calibration. The internal uncertainty was estimated as the standard deviation of the 256 independent measurements. As for MIPS, the total uncertainty in the physical flux density was constructed by adding the internal and absolute uncertainties in quadrature — absolute flux calibration uncertainties of 10 % in all three bands. Low-resolution ($`R=70120`$) spectra were obtained with the Infrared Spectrograph (IRS: Houck et al. 2004) from $`7.438`$µm, using “high-accuracy” blue peak-up to place the source in the spectrograph slit. Integration times per exposure were 6 sec over the Short-Low wavelength range (7.4 $``$ 14.5 $`\mu `$m), and 14 sec over the Long-Low wavelength range (14.0 $``$ 38.0 $`\mu `$m). The spectra for wavelengths beyond $`35\mu `$m are not reliable (Houck et al. 2004), so have been omitted. The spectra were obtained at two nod positions in staring mode for averaging and estimating the noise. The BCDs resulting from the SSC pipeline S10.5.0 were further processed within the SMART software package (Higdon et al. 2004). We used the droopres data products, which are intermediate products produced before stray-light and flat-field corrections have been applied. Spectra were extracted assuming point source profiles with a fiducial width of 5 $``$ 6 pixels in the center of the orders. Residual “sky” emission was subtracted using adjacent pixels. Random errors were calculated from the difference between the two independent spectra, then added in quadrature with an estimated 14 % uncertainty in absolute flux calibration. A more detailed discussion of the IRS data reduction and extraction is presented in Bouwman et al. (2005). ## 3 SPECTRAL ENERGY DISTRIBUTIONS OF FIVE EXO-KUIPER BELTS As part of the ongoing analysis of the FEPS database, we have looked for stars that are detected at 70 $`\mu `$m but show no significant excesses above the stellar photospheres for $`\lambda 33`$µm. This criterion restricts the equilibrium temperature and inner radius of a significant dust component to $`T<70`$ K and $`R_{\mathrm{in}}`$ 10 AU for dust in thermal equilibrium around a star with 1 L. Three other parallel investigations based on FEPS observations will concentrate on stars that have: 1) young optically-thick disks detected by IRAC, which focus on determining the evolutionary time scale of the inner disks (Silverstone et al. 2005); 2) strong emission at 24 µm but weak at 70 µm indicating material at temperatures approaching those of our solar system’s asteroid belt and terrestrial planets (Hines et al. 2005); and 3) infrared excess indicating a broad range of temperatures characteristic of both terrestrial zones and Kuiper Belts (Bouwman et al. 2005). The focus of this study is on stars that are dominated by excess infrared emission suggestive of Kuiper Belt like disks. Five stars were selected from the FEPS data obtained between May 28, 2004 (MIPS campaign 8) and Sep. 25, 2004 (MIPS campaign 13), that have $``$ 3 $`\sigma `$ excess at 70 $`\mu `$m, but are consistent with photospheric emission ($`<`$ 3 $`\sigma `$ excess) $``$ 33 $`\mu `$m. Three of the objects are newly discovered debris systems (HD 6963, HD 145229, & HD 206374), while two of the objects were previously identified as having an infrared excess (HD 8907 & HD 122652; respectively, Silverstone 2000, Zuckerman & Song 2004) from IRAS observations. We note that one object, HD 8907, is one of 10 stars with known excess sources, which was included in the FEPS targets for the purposes of a gas emission line search. A sixth star, HD 13974, was also detected at 70 $`\mu `$m with high signal-to-noise, but in this case the 70 $`\mu `$m flux density is consistent with emission from the photosphere to within 1 $`\sigma `$ uncertainty. Spitzer photometry for the five stars is presented in Table 2, and the SEDs are shown in Figures 1 and 2. Mid-infrared spectra are also displayed for all of the objects. We use the “average weighted” wavelengths of IRAC and MIPS bands as suggested by the Spitzer Observer’s Manual (SOM). Color corrections have not been applied. The magnitude of these corrections are expected to be smaller than the present absolute calibration uncertainties (see the IRAC and MIPS data handbooks<sup>1</sup><sup>1</sup>1http://ssc.spitzer.caltech.edu/). The expected photospheric emission for each star was determined by fitting Kurucz model atmospheres with convective overshoot to published optical photometry including, if available, Johnson $`BV`$, Strömgren $`vby`$, Tycho $`BV`$, Hipparcos $`H_p`$, Cousins $`RI`$ and 2MASS $`JHK_s`$ measurements. Predicted magnitudes in each filter were computed by multiplying the Kurucz model with the combined system response of the filter, the atmospheric transmission (for ground-based observations), and the spectral response function of the detector as outlined in a series of papers by Cohen et al. (2003a, b, and references therein). The Johnson-$`U`$ and Strömgren-$`u`$ filters were excluded from the model fits because the observed photometry shows large deviations from the model values explained in part by chromospheric activity common in young late-type stars. The best-fit Kurucz model was computed in a least squares sense with the effective temperature and normalization constant as free parameters. The metallicity was fixed at \[Fe/H\] = 0.0 and the surface gravity at log g cm s<sup>-2</sup> = 4.5. In addition, each of the stars is within the Local Bubble (see, e.g., Welsh et al. 1998), and the visual extinction was fixed at $`A_V=`$0 mag. The stellar parameters used in the Kurucz models, and additional properties of the stars including their distance and estimated ages, are listed in Table 3. The deviation from the stellar photospheres at $`\lambda 70\mu `$m is readily apparent for HD 6963, HD 8907, HD 122652, HD 145229, and HD 206374 (Fig. 1). The 70 µm emission from HD 13974 is consistent with the stellar photosphere (Fig. 2) within the uncertainty of 70 $`\mu `$m flux. Figure 3 shows the significance of the 70 $`\mu `$m excesses relative to the Kurucz model atmospheres for the five excess stars compared with the distribution for the other stars in the FEPS V1.1 data release. Figure 4 also illustrates the excess in a color-color space that is sensitive to the color temperature of the excess (24$`\mu `$m/Ks vs. 70$`\mu `$m/Ks). Large filled circles in Figure 4 are sources with 70 µm detections, and small open circles with arrows are 1 $`\sigma `$ upper limits for non-detections. This plot has the virtue of also capturing the lack of strong 24 $`\mu `$m excesses in these targets. Note that the 24 $`\mu `$m excess of HD 12039 is well separated from the other targets; this star is unique within the FEPS program in exhibiting only warm debris dust, and is discussed in detail in Hines et al. (2005). The location of HD 13974 in this plot agrees well with 1 $`\sigma `$ upper limits of 70 µm non-detections. In the next section we use the observed SEDs to constrain the properties of the detected debris systems. For HD 13974 the 70 µm detection is consistent the photosphere, placing stringent limits on any KB-like excess. For those objects that were not detected at 70$`\mu `$m, we are still able to place upper limits on the presence of a debris system. ## 4 DEBRIS DISK MODELS We model our debris disks assuming optically thin dust in thermal equilibrium with the stellar radiation field. In this case, the temperature of a dust grain with a given chemical composition and grain size depends on the radial distance to the central star only. The fact that the SEDs do not significantly deviate from the photosphere for $`\lambda 33`$ µm places the maximum temperature for the detected dust grains at $`100`$ K. This suggests a minimum equilibrium distance from the stars $`10`$AU, assuming grains similar in size or larger than those found in our own zodiacal dust cloud (radius $`a10100`$ $`\mu `$m, Reach et al. 2003). The dust grains are subject to radiation pressure and Poynting-Robertson (P-R) drag, and the action of both mechanisms places limits on the time that the dust remains in the system. Very small grains will be ejected from the system by radiation pressure, while larger grains will suffer P-R drag and spiral toward the parent star. A detailed prescription for computing “blowout” sizes and P-R drag timescales is given in Burns et al. (1979) for the solar system, and are derived for a range of extrasolar systems by Artymowicz (1988) and Backman & Paresce (1993). Particles smaller than $`1\mu `$m have a blown out time $`<`$10<sup>2</sup> yrs. We consider grains to be no longer important contributors to the SED when they travel a factor of 4 farther from the star than their position of creation in the main dust belt. At these distances, the grain is 0.5 $`\times `$ initial temperature, and emits at a twice the wavelength. We consider such a grain to no longer be contributing to the main excess SED, and in effect has left the system. Grains blown out from points of origins a few $`\times `$ 10 AU from a solar type star will leave a system by this definition in at most only a few$`\times `$10<sup>2</sup> years. Particles larger than $``$ 1 µm will be subject to slow P-R inward drift and will be destroyed on timescales of 10$`{}_{}{}^{6}10^7`$ yrs for particle sizes of $``$1 $``$ 10 µm (timescale linearly proportional to particle size) starting from distances of order $`r`$ $`10`$ AU (timescale proportional to $`r^2`$). This time scale is very short compared to the age of the systems (see Table 3), which suggest that the debris in these systems is being replenished by a parent population of objects, a putative Kuiper Belt. From these arguments, we can build simple models of the debris systems. Below we develop models for five of the objects that have limited photometry beyond 24 $`\mu `$m. For HD 8907, where we have more photometric constraints on the SED, we explore a more comprehensive model. ### 4.1 Simple Blackbody Grain Models Blackbody models of the detected debris disks were based on excess color temperature ($`T_\mathrm{c}`$), estimated from the two shortest broadband wavelengths with $``$ 1 $`\sigma `$ excess (including the synthetic band at 33 µm using IRS long-low data) — Blackbody grains by definition absorb and emit radiation efficiently at all relevant wavelengths. For each object, the solid angle, $`\mathrm{\Omega }`$, subtended by the emitting material as seen from the central star was calculated from $`T_\mathrm{c}`$ and excess $`F_\nu `$ using the Planck formula. This value for $`\mathrm{\Omega }`$ in turn was converted to: (1) total grain emitting cross-section area, $`A_\mathrm{x}`$, using the system distance, (2) grain luminosity using $`A_\mathrm{x}`$ and $`T_\mathrm{c}`$ in the Stefan-Boltzmann formula for blackbody luminosity, and (3) grain mass assuming a constant per-grain mass density <sup>2</sup><sup>2</sup>2A density of 2.5 g cm<sup>-3</sup> was chosen to represent generic silicate fragments of asteroid-like planetesimals, c.f., mean value of densities determined for 10 asteroids $``$2.4 g cm<sup>-3</sup>; http://aa.usno.navy.mil/hilton/asteroid\_masses.htm. of 2.5 g cm<sup>-3</sup> for solid “astronomical silicate” (Draine & Lee 1984) and grain radius 10 µm for efficient emission at 70 µm. The 10 µm radius grain is the smallest size that acts as a blackbody at 70 µm. The uncertainty in blackbody temperature estimate ($`T_\mathrm{c}`$) is about 10%. In all cases, the flux densities between 33 $`\mu `$m and 70 $`\mu `$m are increasing, suggesting that we are observing the Wien side of the dust excess SED, thus $`T_\mathrm{c}`$ is interpreted as an estimate of the maximum dust temperature. The corresponding minimum distance from the star of the emitting material $`R_{\mathrm{in}}`$ was found from the estimated maximum $`T_\mathrm{c}`$ and the relevant $`L_{}`$ using the formula for blackbody grains in Backman and Paresce (1993). The relationship between grain temperature, position ($`R_{\mathrm{AU}}`$), and primary star luminosity (equation 3 of Backman & Paresce 1993) for grains larger than $`a10\mu `$m, that emit efficiently at $`70\mu `$m, and have negligible albedo, is $$T_\mathrm{c}=278L_{}^{\frac{1}{4}}R_{\mathrm{AU}}^{\frac{1}{2}}\mathrm{K}$$ (1) where $`L_{}`$ is in unit of $`L_{\mathrm{}}`$, 3.9$`\times `$10<sup>33</sup> erg s<sup>-1</sup>. The model fit parameters and masses are presented in Table 4, and Figure 1 shows simple model SEDs. The blackbody models also allow estimates of the mass in the radiating grains. Because (1) these calculations made use of a maximum dust temperature and corresponding inner radius ($`R_{\mathrm{in}}`$), also (2) lack of photometric measurements or even useful upper limits beyond the apparent peak of excess emission (70 $`\mu `$m) prevents useful characterization of the dust outer boundary radius, and finally (3) grains larger than the minimum size assumed here would radiate the same blackbody SED but would have a higher ratio of mass to surface area, the values of $`L_{\mathrm{IR}}/L_{}`$ and $`M_\mathrm{d}`$ in Table 4 characterize the portion of the SED observable out to 70 µm; future observations at longer wavelengths could conceivably indicate more luminosity and dust mass. We note that debris disk models comprised of grains slightly larger than the blowout size ($``$ 1 µm) cannot be ruled out based on simple fits to the SED. Such a model (based on a temperature vs. position function given in equation 5 of Backman and Paresce 1993) requires an $`R_{\mathrm{in}}`$ of $``$ 270 AU for a solar luminosity star and a dust mass $``$ 3 larger than the lower limit given above. We comment on the likelihood of this model in section 5.4. ### 4.2 Detailed Modeling of HD 8907 The simple black body grain model yields $`T_\mathrm{c}=48`$ K, the inner radius $`R_{\mathrm{in}}`$ = 48 AU, a dust to stellar luminosity ratio (log $`L_{\mathrm{IR}}/L_{}`$) $`=3.64`$, and a dust mass log ($`M_\mathrm{d}/M_{\mathrm{}}`$) $`=3.08`$ for 10 $`\mu `$m radius blackbody grains with grain density 2.5 g cm<sup>-3</sup>. The detection of significant excess emission from HD 8907 at several wavelengths allows us to conduct a more thorough analysis of that system. Therefore we modeled the SED of HD 8907 using the dust disk models of Wolf & Hillenbrand (2003, WH03; 2005). We used the Levenberg-Marquardt algorithm to solve the least-squares problem, giving a best-fit dust disk model (Marquardt 1963; Markwardt 2003 <sup>3</sup><sup>3</sup>3http://cow.physics.wisc.edu/$``$craigm/idl/fitting.html). A detailed description of the fitting algorithm will be given in Rodmann et al. (2005). Due to lack of mineralogical features in the IRS spectrum, we assumed the optical properties of astronomical silicates (Draine & Lee 1984; Laor & Draine 1993; Weingartner & Draine 2001). We assumed the density profile to be $`n(r)r^1`$, corresponding to a disk with a constant surface density $`\mathrm{\Sigma }(r)r^0`$. The power law exponent of the grain size distribution $`n(a)a^s`$ was set to the canonical value $`s=3.5`$, characteristic for a size distribution initially produced by a collisional cascade (Dohnanyi 1969; Tanaka et al. 1996). For the maximum grain size ($`a_{\mathrm{max}}`$) and the $`R_{\mathrm{out}}`$, we assumed arbitrary values of $`a_{\mathrm{max}}=1`$mm and $`R_{\mathrm{out}}=100`$AU. Grains larger than 1 mm contribute little to the infrared emission of the disks, and the outer radius is poorly constrained in the absence of sub-millimeter/millimeter data. We used an upper limit at $`\lambda =3.1`$mm (Carpenter et al. 2005) to constrain flux of the disk models at mm wavelengths. We then fit for three parameters simultaneously in a $`\chi ^2`$ sense: $`R_{\mathrm{in}}`$, $`a_{\mathrm{min}}`$, and M<sub>dust</sub> in grains smaller than 1 mm. The first two parameters influence the shape of the SED and the wavelength at which the disk emission begins to depart from the stellar photosphere; the latter only scales the dust reemission to match the peak of the infrared emission. Figure 5 shows the SED of HD 8907 with the WH03 model result using astronomical silicate and grain radii range from 6 µm $``$ 1 mm. The result suggests $`R_{\mathrm{in}}`$ $``$ 42.5 AU, with dust mass of 1.7$`\times 10^2`$ $`M_{\mathrm{}}`$. The larger dust mass compared to the mass estimate from the simple blackbody model is from the inclusion of larger grains, which contribute the bulk of the mass even though they do not contribute to the bulk of the infrared emission. The 1 $`\sigma `$ uncertainty level for each parameters used in our modeling are: $`\delta `$($`R_{\mathrm{in}}`$) $``$ 20 AU (50%); $`\delta `$($`a_{\mathrm{min}}`$) $``$ 3 µm(50%); and $`\delta `$($`M_\mathrm{d}`$) $``$ 3.3 $`\times `$ 10<sup>-3</sup> $`M_{\mathrm{}}`$ (20%). These parameter uncertainties are estimated from the covariance matrix (inverse of the $`\chi ^2`$ curvature function) by the Levenberg$``$Marquardt algorithm. ### 4.3 On the Presence of Warm Dust Mass Within $`R_{in}`$ The selection of debris disks for this investigation is purposely biased against objects with debris at temperatures warmer than $`100`$ K. However, in our own solar system we observe a warm component associated with debris from the asteroid belt. Such a warm component, albeit faint, has also been verified in other debris systems, e.g., around an A-star HR 4796A (e.g., Koerner et al. 1998), and recently with Spitzer imaging of Fomalhaut (Stapelfeldt et al. 2004). We do not have the luxury of high resolution imaging of the five Sun-like stars presented herein, so we have to rely on the SEDs to place upper limits on the warm dust components. Table 5 presents upper limits to the amount of warm dust located interior to the outer disk inner boundaries ($`R_{\mathrm{in}}`$) in each system. These were evaluated by assuming that dust in these locations would be drifting in from the denser outer source zones via P-R radiation drag and thus would extend with constant surface density from $`R_{\mathrm{in}}`$ inward to a radius correspond to the vaporization temperature of silicates at $`T=`$ 1500 K at $``$0.1 AU from the star. For each system a surface density ($`\mathrm{\Sigma }`$) upper limit was calculated for dust between sublimation radius, $`R_{\mathrm{sub}}`$($`T=`$1500 K) and $`R_{\mathrm{in}}`$ of the outer disk model such that the summed SEDs of inner dust model plus outer dust model plus stellar photosphere was less than the upper limit to the observed SED, i.e., observations plus 1$`\sigma `$ uncertainty including calibration, as defined above. This was then converted to dust mass limits assuming $`a`$10$`\mu `$m grains with material density of 2.5 g cm<sup>-3</sup> as for the blackbody models of the cool outer disks. These masses are of order 10<sup>-6</sup> M, the mass of a single asteroid only a few hundred km in diameter, and are generally 2$``$3 orders of magnitude below the lower limits for the masses we derive for the outer cool dust. Note that our Spitzer observations of HD 13974, which shows photospheric emission up to 70 $`\mu `$m, yield an upper limit on warm dust surface density equivalent to only about 20 times that of the zodiacal cloud in our solar system (i.e., 20 ”zodis”). ## 5 DISCUSSION A primary goal of the FEPS project is to place our solar system in context with other debris disk systems (Meyer et al. 2005). To date, most disk systems identified with IRAS and ISO have been associated with either stars more luminous than the Sun (typically A-stars), or very young systems ($`t_{\mathrm{age}}`$ 30 Myrs). A recent census of Sun-like stars with KB-like disks that have been identified with IRAS and ISO has been presented by Decin et al. (2003) and Zuckerman & Song (2004), but these studies were only able to identify the brightest systems and did not cover a large range of ages that encompass stars approaching the age of our Sun. The FEPS sample is a coherent effort to understand the debris systems around Sun-like stars over a large range of ages from 3 Myrs – 3 Gyrs. We are beginning to assemble a sufficiently large sample of objects over a range of ages to start to answer fundamental questions about the evolution of dust disks surrounding Sun-like systems with Spitzer data. ### 5.1 Age Determination The age determination of all of the stars in the full FEPS sample is discussed in Hillenbrand et al. (2005). Age bins given for the five stars (Table 3) in this study are inferred from the level of chromospheric and coronal activity. The chromospheric activity is indicated by CaII H&K emission. Values of the fractional luminosity emitted in the H&K lines, $`R_{HK}^{}`$, were derived by Soderblom (2000, private communication). Here we adopt the calibration of Donahue (1993) between log $`R_{HK}^{}`$ and age. Coronal activity, on the other hand, is indicated by X-rays. The fractional X-ray luminosity, $`R_\mathrm{X}=L_\mathrm{X}/L_{\mathrm{bol}}`$, in general tracks $`R_{\mathrm{HK}}^{}`$; for HD 8907 in particular, relative youth is indicated by both activity indicators. We use coarse age bins rather than specific ages due to the uncertainty in the age estimates. ### 5.2 Solar System Evolutionary Model The total mass of the KB in the solar system is highly uncertain, however some recent studies (e.g., Luu & Jewitt 2003, Teplitz et al. 1999) suggest its total mass is in the range $``$0.01 $``$ 1.0 M, mostly in large planetesimal bodies. From COBE observations at wavelengths 150 $`\mu `$m and 240 $`\mu `$m, Backman et al. (1995) estimated upper limits on KB dust luminosity of $`L_{\mathrm{IR}}/L_{}`$ 10<sup>-6</sup>. This corresponds to a KB infrared dust mass limit of $`M_{\mathrm{d},\mathrm{KB}}`$10$`{}_{}{}^{5}M_{\mathrm{}}^{}`$ assuming grains with $`a=`$ 10 $`\mu `$m and $`\rho `$ = 2.5 g cm<sup>-3</sup> as in the present paper’s models of exo-KBs. For comparison, Moro-Martín & Malhotra (2003) estimate a KB mass of about 4$`\times `$10<sup>-6</sup>M in particles with sizes between 2.4 $`\mu `$m and 160 $`\mu `$m based on flux of KB dust detected drifting toward Jupiter’s orbit (Landgraf et al. 2002). We see that the systems considered in this paper contain more massive and luminous dust ensembles than does the KB. But what would our own KB have looked like at the younger evolutionary stages represented by these five objects? Backman et al. (2005) have devised a simple evolutionary model of the KB based on the model of the current KB in Backman et al. (1995) that can be used to estimate what our system would have looked like at the ages, distances and stellar luminosities of our targets. The model assumes the planetesimal population in the KB of these solar systems extends from $`R_{\mathrm{in}}`$ $`=`$ 40 AU to $`R_{\mathrm{out}}`$ $`=`$ 50 AU at the ages of our targets, after migration of the outer planets and KB had been completed (Malhotra 1993; Levison et al. 2004), and the influence by Neptune (at 30 AU) had eroded the belt substantially inward of the 3:2 resonance at R $`=40`$ AU. The model also assumes that the starting mass of the KB was 10 M, the minimum necessary to yield a high enough density to build the observed large KB bodies accretionally in the allowed time span (Stern & Colwell, 1997). The model evolution is not very sensitive to the starting mass because a purely collisional system evolves asymptotically toward a state in which the collision timescale is approximately equal to the age of the system. For example, a wide range of starting masses evolve in this model to a KB with a mass of about 0.5 M at age 4.5 Gyr. This means that, in a dissipative system such as this, evidence of the original state is mostly erased and cannot be inferred by modeling backwards. A parent body size of 10 km and maximum fragment size of 5 km were parameters tuned to yield a KB dust population at an age 4.5 Gyr consistent with present limits on KB dust emission (Backman et al. 1995; Teplitz et al. 1999) and inference of the KB dust production rate from dust impact rates in the outer solar system (Landgraf et al. 2002). The dust distribution and thermal emission were calculated for 30 logarithmically spaced size bins between 1 $`\mu `$m and 1 mm diameter. An equilibrium spatial distribution of dust was calculated, as in Backman et al. (1995), balancing the dust between production in planetesimal collisions, inward drift of dust via P-R radiation drag, and destruction by mutual collisions of dust yielding fragments smaller than the blowout size. The dust distribution was modeled both within the planetesimal zone at 40 $``$ 50 AU and in a P-R induced “zodiacal” dust cloud extending inward and assumed completely truncated at $`R_{\mathrm{out}}`$ $`=`$ 30 AU. Results of the KB evolutionary model runs are presented in Table 6, where columns 3 and 4 are the predicted and observed 70 $`\mu `$m excess fluxes of each source. Remarkably, although the solar system evolutionary model was tuned only to represent our KB, the results are within a factor of only 2 $``$ 3 from the targets’ observed 70 $`\mu `$m excesses, except for the HD 13974. The solar system KB evolution model predicts a 70 µm excess of $``$ 21 mJy for the Sun observed from $`d=`$ 30 pc at an age of $`t=`$4.6 $`\times 10^9`$ yrs (on top of 5 mJy photosphere). This predicted excess flux is less than any of the 5 detected excess sources, e.g., about 30 % of the dust flux at 70 µm observed for HD 145229 with 0.96 L observed at 33 pc. Therefore the solar system present-day KB model has about 30 % of the dust mass of the HD 145229 cold dust. These calculations indicate the possibility that the target systems represent snap shots of the history of our own solar system suggesting that systems like our own might be common among G stars in the galactic disk. ### 5.3 HD 13974 HD 13974 is a short period binary system (period = 10 days, Duquennoy & Mayor 1988) with a G0V primary (Duquennoy & Mayor 1988) and a companion with spectral type between G9V (Duquennoy & Mayor 1988) and K4V (Hummel et al. 1995). HD 13974 at 11 pc is the only star in our FEPS sample to date with a 70 $`\mu `$m detection that is consistent with “bare” photospheres within the calibration and Kurucz model uncertainties (Figure 2). A two-component Kurucz model was fitted to the observed SED with derived temperatures of 6215 K and 4493 K for the primary and secondary components respectively. The semi-major axis of these stars are $`a_1sini`$ $`=`$ 1.45 $`\pm `$ 0.03 $`\times `$ 10<sup>6</sup> km and $`a_2sini`$ $`=`$ 1.62 $`\pm `$ 0.04 $`\times `$ 10<sup>6</sup> km (Duquennoy & Mayor 1988). Such a close ($``$ 1 AU) binary star system should not have a dramatic effect on the evolution of the circumbinary disk at distances greater than roughly twice the semi-major axis (e.g., Jensen et al. 1996; Artymowicz & Lubow 1994 ). Adopting the blackbody grain model we have used the 1$`\sigma `$ uncertainty in the observed flux density at 70 $`\mu `$m compared with the predicted Kurucz model flux to place an upper limit on the mass of dust log ($`M_\mathrm{d}/M_{\mathrm{}}`$)$`=5.1`$ at $`T_\mathrm{c}=`$ 55 K. The lower limit on $`R_{\mathrm{in}}`$ is about 28 AU. The fractional dust luminosity log ($`L_{\mathrm{IR}}/L_{}`$) $`<5.2`$ of this system is interesting, as it is similar to the estimates for the current solar system Kuiper Belt. The solar system KB evolutionary model predicts 70 $`\mu `$m excess flux density of HD 13974 at its age (1 $``$ 3 Gyr) to be about 300 mJy, while the observed excess flux density limit is $`<`$ 20 mJy (Table 6). This system may not contain KB bodies, or perhaps there is not a perturbing planet like Neptune to stir up the system and cause collisional cascades. ### 5.4 Possible Planetary Architectures in the Five IR-Excess Stars For the five systems in which 70 $`\mu `$m excesses have been detected, the comparison between the disk mass estimate ($`M_\mathrm{d}`$ in Table 4), and the upper limit to the amount of warm dust ($`M_{\mathrm{d},\mathrm{warm}}`$ in Table 5) located interior to the inner boundary of the dust ($`R_{\mathrm{in}}`$ in Table 4), indicates that the dust depletion inside $`R_{\mathrm{in}}`$ is significant. Because grains not directly ejected by radiation pressure tend to spiral toward the star due to P-R drag, a central depletion would be filled in by dust on P-R timescales (Table 5) much shorter than the ages of these systems, unless some other mechanism intervenes to eliminate grains. We argue that ice sublimation and grain “blowout” both fail to explain the $`R_{\mathrm{in}}`$ locations in these systems, and that a possible explanation for the dust spatial distribution is a sizable planet in each system limiting the inward drift of the grains by ejecting them out of the system via gravitational scattering. An upper limit of the location ($`R_{\mathrm{sub}}`$) at which grains will sublimate can be calculated by assuming that they are composed of water ice with sublimation temperature $`T_{\mathrm{sub}}=`$100 K. The maximum grain temperatures ($`T_\mathrm{c}`$) observed in these systems presented in Table 4 are all well below the ice sublimation temperature. The presence of grains large enough to radiate efficiently at 70 $`\mu `$m implies that their temperature will scale as $`T_\mathrm{c}`$ $`R^{0.5}`$ (Backman & Paresce, 1993). The radius of ice sublimation $`R_{\mathrm{sub}}`$ will be $`1/41/3`$ of $`R_{\mathrm{in}}`$ for $`T_\mathrm{c}`$ in the range $`5070`$ K. and thus sublimation is not a likely explanation for the inner depletion of the disk. The uncertainties in the temperature estimate are about 10% for individual sources. Can these central zones that are relatively free of dust be explained as the result of dust grains in the outer rings being controlled by mutual collisions rather than P-R drag? Thus the inner edge of the dust distribution might be explained without invoking the existence of a planet to consume inward drifting grains. In their study of the age dependence of Vega-like excesses, Dominik & Decin (2003) argued that in dense enough disks, the collision time scale is much shorter than the P-R drag time scale. If grain collisions actually dominate P-R drag as the dominant dynamical process then grains could be ground to the blowout size, and would not drift radially. In that case, the inner edge to the dust distribution could represent the inner edge of the source bodies. The dust collision timescales ($`t_{\mathrm{coll}}`$) for the disks can be estimated as $`P_{\mathrm{orb}}`$/9$`\mathrm{\Sigma }`$, where $`P_{\mathrm{orb}}`$ is the orbital period at $`R_{\mathrm{in}}`$ and $`\mathrm{\Sigma }`$ is the fractional surface density of the disks that is of order the value of $`L_{\mathrm{IR}}/L_{}`$ (Backman & Paresce 1993; Backman 2004). The timescales for particles to drift via P-R drag over distances equal to 10% of the $`R_{\mathrm{in}}`$ values in table 4 (i.e., an estimate of the time for grains to move from populated to relatively unpopulated regions) are about log $`t_{\mathrm{P}\mathrm{R}}`$ 5.5 to 6.0 years. The grain-grain collision timescales for the same systems are about log $`t_{\mathrm{coll}}`$ 5.0 to 5.5 years. For ratios of this sort between $`t_{\mathrm{P}\mathrm{R}}`$ and $`t_{\mathrm{coll}}`$, a substantial fraction of grains would survive enough collision lifetimes to be able to drift from the disk into a central void. Note also that without knowing more about the disks (especially, their vertical extents) we do not know if grain-grain collisions would be at high enough speeds to destroy the grains. Thus we conclude that P-R drag is important in controlling the structures of these systems, especially in that a planetary barrier would be a plausible explanation for a central depletion. Could these cold disks arise from grains generated recently from planetesimals located near $`R_{\mathrm{in}}`$ and then ejected by radiation pressure from the system, again without needing a planet to explain the inner boundary of the disk? We note that these debris belts can be self-stirred without needing a planet to explain the collisional cascade itself (e.g., also Kenyon & Bromley 2004). Recent Spitzer observations of Vega at 24, 70, and 160 $`\mu `$m show that the Vega debris disk has an inner boundary at $``$ 86 AU, and extends to distances much larger than those observed at submm and mm wavelengths (Su et al. 2005). The authors suggest that the grains in the Vega far-IR disk are produced close to the inner disk boundary and are unbound, flowing away from the star. In this scenario, the inner gap in the disk would simply represent a limit to the location of the parent bodies, rather than the location of a perturbing planet. However, Vega’s luminosity to mass ratio (L/M) is 24 in solar units, which makes even relatively large grains (like the 18 µm grains) to be unbound, with $`\beta >`$ 0.5, if the grain’s porosity is $`>`$0.56. ($`\beta `$ is the ratio of the radiation pressure force to the gravitational force). In contrast, for the G star systems we consider here with L/M $``$ 1, the 10 µm silicate grains that dominate the dust emission at 70 µm have $`\beta =`$ 0.02. Therefore, we argue that grains in the five disks discussed here are bound and as such will drift toward the star due to P-R drag. One explanation for the inner empty zones is that one or more massive planets dynamically deplete the dust generated by an outer belt of planetesimals (e.g., Liou & Zook 1999, Moro-Martín & Malhotra 2003, 2005). If we assume a single planet in a circular orbit about each of the five stars, then this mechanism can account for the large dust depletion factors inferred for the inner parts of these dust disks if the planet is $`>`$ 1 M<sub>Jupiter</sub> (Moro-Martín and Malhotra 2005). The inner edge of the dust disk would be located between 0.8 $`\times a_{\mathrm{pl}}`$ and 1.25 $`\times a_{\mathrm{pl}}`$, where $`a_{\mathrm{pl}}`$ is the semi-major axis of the planet. Therefore, a possible explanation is that the dust disk surrounding each of these five systems harbors a planet with an orbital radius of approximately 10 $``$ 20 AU, together with a belt of dust-production planetesimals exterior to the planet’s orbit, and no asteroid belt with significant dust production interior to the planet’s orbit. As noted in §4.1, the data are also consistent with a debris disk comprised of 1 µm sized grains, inconsistent with the assumptions required for blackbody emission discussed above, yet larger than the blowout size. In such a model, the timescale for mutual collisions is more than 100 times shorter than the timescale for P–R drag to significantly affect the orbital radius. This suggests that a ring such as inferred from the model could be maintained through mutual collisions with dust removal dominated by radiation pressure blowout of the smallest particles $`<`$ 1 $`\mu `$m on short timescales. However, this requires the formation of a sun–like star lacking an inner disk, with a remnant planetesimal belt capable of generating dust through mutual collisions only outside of 200 AU with dust only at the smallest possible stable particle size. While the presence of a low mass companion could in principle explain such a large gap in the circumstellar environment, no wide stellar companions are known in these systems, and we do not expect many brown dwarfs (due to the observed lack of such systems) or giant planets (given current formation scenarios) at these radii. We believe that models requiring larger grains ($``$ 10 µm) at smaller radii ($``$ 30 AU) with the possible presence of a planet (or planetary system) interior to 30 AU are favored as they are consistent with the structure and extent of circumstellar disks observed around pre–main sequence stars during the epochs when planetary systems are thought to form. We hope to distinguish between these two hypotheses in the future through direct imaging of the outer edge of the disks in scattered light and placing strong constraints on the presence of very low mass companions at large separations (e.g. with NICMOS coronagraphy on HST and ground–based AO observations). ## 6 SUMMARY We have presented Spitzer observations of five Sun-like stars that possess Kuiper Belt-like debris systems (HD 6963, HD 8907, HD 122652, HD 145229, and HD 206374). Of these two were previously suggested to have debris disks (HD 8907 and HD 122652), and three are newly discovered. The fractional luminosity ($`L_{\mathrm{IR}}/L_{}`$) and ages of these newly identified disks illustrates the potential of the FEPS program to measure fainter debris systems, which will allow a more detailed census of their nature and evolution. We summarize our results: 1. The five excess sources have SEDs that are consistent with photospheric models out to 33$`\mu `$m, but show clear excesses at 70 $`\mu `$m, which was the selection criterion. We find that these stars are all ”old” (four sources are in our 1 $``$ 3 Gyr age bin, while one, HD 145229, is in the 0.3 $``$ 1 Gyr age bin). 2. As seen in Figure 3, the improved sensitivity of Spitzer allows us to detect debris disk systems that are much fainter than those detected by IRAS and ISO. The overall impression is that KB-like systems detectable by Spitzer and considered in this paper are less massive and more distant than systems detected with IRAS and ISO. 3. Another star, HD 13974, has a MIPS 70 $`\mu `$m flux consistent with photospheric emission within 1 $`\sigma `$ total uncertainty. The upper limit of log($`L_{\mathrm{IR}}/L_{}`$) is $`<5.2`$, similar to that of inferred for the solar systems’ KB. 4. Simple blackbody grain modeling of our 5 excess SEDs yielded log($`L_{\mathrm{IR}}/L_{}`$) $`4.53.5`$, color temperatures between 55 $``$ 58 K, and inner radii of outer disks between 18 and 46 AU. 5. A solar system KB evolution model predicts Spitzer 70 $`\mu `$m fluxes (Table 6) from hypothetical planetesimal assemblages around our target stars that are within factors of $``$2 $``$ 3 of the observed fluxes. We infer that these systems have outer remnant planetesimal belts that are consistent in scale and starting masses to our Kuiper Belt. 6. The absence of a disk around the $`1`$ Gyr old star HD 13974 suggests that either this object does not contain the parent bodies that produce infrared-emitting debris, or perhaps the debris has been cleared out already. 7. We placed upper limits on warm dust masses interior to $`R_{\mathrm{in}}`$ for each of these systems, and showed that the depletion of the disk $`<`$$`R_{\mathrm{in}}`$ is significant. We commented on several possible causes for $`R_{\mathrm{in}}`$. We speculate that the $`R_{\mathrm{in}}`$ of exo-KBs presented in this study could be explained by the existence of one or more Jupiter mass planets at 10 $``$ 20 AU from each star. We thank to the rest of the FEPS team members and Spitzer Science Center help desk. We have used the SIMBAD database. This work is based \[in part\] on observations made with the Spitzer Space Telescope, which is operated by the Jet Propulsion Laboratory, California Institute of Technology under NASA contract 1407. FEPS is pleased to acknowledge support through NASA contracts 1224768, 1224634, and 1224566 administered through JPL. S.W. was supported by the German Research Foundation (DFG) through the Emmy Noether grant WO 857/2-1. EEM is supported by a Clay Fellowship from the Smithsonian Astrophysical Observatory. MPIA team is supported through the European Planet Network.
warning/0506/cond-mat0506453.html	ar5iv	text	# Self-Organized Ordering of Nanostructures Produced by Ion-Beam Sputtering ## Model.– During IBS, the bombarding ions penetrate the target and induce complex collision cascades in the bulk. In semiconductor substrates like those studied in facsko\_science ; gago ; frost , these cascades amorphize the near-surface layer. Sputtering events take place when surface atoms receive enough energy and momentum to break their bonds and leave the target. We will assume that only a fraction of those atoms are redeposited at the surface. Adatoms are moreover available to relaxation mechanisms such as surface diffusion, that can be thermally activated, or else be induced by the mentioned change in the local viscosity of the material close to the surface viscousflow . In the spirit of the so-called hydrodynamic theory of ripples in aeolian sand dunes sand , we define two coupled fields, namely, $`R(𝐱,t)`$ and $`h(𝐱,t)`$, where $`𝐱=(x,y)`$. The first one represents the fraction of surface atoms that are not sputtered away but, rather, remain mobile along the target surface. Analogously, $`h`$ measures the height of the surface neglecting the contribution from the fraction of mobile atoms $`R`$. Time evolutions of $`R`$ and $`h`$ are coupled through reaction and transport mechanisms aste . Thus, $`_th`$ $`=`$ $`\mathrm{\Gamma }_{ex}+\mathrm{\Gamma }_{ad},`$ (1) $`_tR`$ $`=`$ $`(1\varphi )\mathrm{\Gamma }_{ex}\mathrm{\Gamma }_{ad}𝐯R𝐉,`$ (2) where $`\mathrm{\Gamma }_{ex}`$ and $`\mathrm{\Gamma }_{ad}`$ are, respectively, the rates of excavation and addition to the surface, $`𝐯`$ is the average velocity of mobile atom, and $`\varphi 0`$ is the fraction of adatoms that detach irreversibly from the surface. Thus, system (1)-(2) does not conserve the amount of material, in marked contrast with typical conditions for aeolian sand dunes dunes\_non\_cons . Here, large redeposition of sputtered atoms corresponds to the small $`\varphi `$ limit, while, in the absence of redeposition, $`\varphi =1`$. Considering that matter transport along the surface is due to diffusion of mobile species, we set $`𝐉=D^2R`$, where $`D`$ is the surface diffusivity. In the absence of bombardment, the concentration of mobile adatoms $`R`$ changes due to thermal nucleation of adatoms from the “immobile state” $`h`$, and subsequent transport along the surface. Assuming nucleation events are more likely in surface protrusions, we have $`\mathrm{\Gamma }_{ad}^{\mathrm{no}\mathrm{er}.}=\tau ^1[RR_{eq}^0(1+\mathrm{\Lambda }\kappa )]`$, analogous of the Gibbs-Thompson relation, $`\kappa `$ being the mean surface curvature and $`\mathrm{\Lambda }`$ the capillary length, assumed isotropic due to amorphization by the ion beam. Here $`\tau `$ is related to the mean time between nucleation events, and $`R_{eq}^0`$ is the mean equilibrium concentration of mobile species for a flat surface. In the presence of bombardment, $`\mathrm{\Gamma }_{ad}^{\mathrm{no}\mathrm{er}.}`$ has to be generalized, to include the contribution of erosion to surface mobility viscousflow . If the ions fall onto the target along the $`x`$ direction, forming angle $`\theta `$ with the normal to the uneroded target, we have, for small slopes tbp ; nota4 , $`\mathrm{\Gamma }_{ex}`$ $`=`$ $`\alpha _0[1+\mu _2(h)^2](1+𝜶_1h+\alpha _2^2h)`$ (3) $``$ $`\alpha _0[\alpha _3(h)^2\alpha _4(_xh)(^2h)]\beta ^2h,`$ $`\mathrm{\Gamma }_{ad}`$ $`=`$ $`\gamma _0[RR_{eq}(1\gamma _2^2h)],`$ (4) where $`R_{eq}`$ and $`\gamma _i`$ generalize parameters in $`\mathrm{\Gamma }_{ad}^{\mathrm{no}\mathrm{er}.}`$ so that $`\gamma _0=\tau ^1+\tau _{ex}^1`$, $`\gamma _2=\mathrm{\Lambda }+\mathrm{\Lambda }_{ex}`$, with $`\tau _{ex}^1`$ and $`\mathrm{\Lambda }_{ex}`$ being analogs of nucleation time and capillary length of erosive origin sand ; viscousflow . Coefficients $`\alpha _i0`$ in (3) are related to geometric correction factors that take into account the local variation of the ion flux with the surface slopes tbp . $`E.g.`$, for oblique incidence, $`\alpha _1,\alpha _4\mathrm{sin}\theta `$, and $`\alpha _3=1/2`$. Likewise, coefficient $`\mu _20`$ is related to the local variation of the sputtering yield with the surface slope new\_note , assumed to have a local minimum for normal incidence, while $`\beta 0`$ measures the efficiency of erosion due to direct impingement of the ions onto surface atoms (knock-on sputtering) reviews\_ibs\_morph ; sigmund . The positive sign of $`\alpha _2`$ implements the physical instability inherent to Sigmund’s theory, by which erosion is more efficient at surface depressions than at surface protrusions sigmund . Actually, the analysis presented below will allow us to relate some of these coefficients with the parameters characterizing Sigmund’s distribution of energy deposition. ## Surface dynamics.– Our continuum model of IBS, formed by Eqs. (1)-(2), (3)-(4), provides a way to introduce systematically all relevant physical mechanisms for IBS, differing from that in aste in a number of features. Rather than considering its full solution, we proceed by deriving an effective equation for the surface height. As in the experiments of references facsko\_science ; frost ; gago , we consider the case of ions bombarding the target at normal incidence ($`\theta =0`$), thus $`\alpha _1=\alpha _4=\|𝐯\|=0`$ in (2), (3) nota . After a transient time of order $`\gamma _0^1`$, Eqs. (1)-(2) have a planar solution $`h_0(t)=\alpha _0\varphi t`$, $`R_0(t)=R_{eq}+(1\varphi )\alpha _0/\gamma _0`$. Perturbing this solution with periodic waves of the form $`h_k=\stackrel{~}{h}_k\mathrm{exp}(\omega _kt+\mathrm{i}𝐤𝐱)`$, and an analogous expression for $`R_k`$, amplification/decay of such perturbations is characterized by the dispersion relation $`\omega _k=R_{eq}\gamma _0\gamma _2\left(ϵ\varphi k^2\gamma _0^1(D+\varphi A)[1ϵ(1\varphi )]k^4\right)`$, with $`ϵ=A/(R_{eq}\gamma _0\gamma _2)`$ and $`A=\alpha _0\alpha _2\beta `$. If $`A>0`$ in $`\omega _k`$, $`i.e.`$ if sputtering is dominated by collision cascades rather than knock-on events, as occurs at low to intermediate energies where Sigmund’s theory is applicable, there is a band of unstable modes that grow exponentially fast, with a linear dispersion relation $`\omega _k`$ of the expected KS type. At this stage, the surface morphology is dominated by a periodic pattern whose wave-vector maximizes $`\omega _k`$. In-plane isotropy under normal incidence implies dependence of $`\omega _k`$ on $`k=\|𝐤\|`$ rather than the full wave-vector $`𝐤`$, thus the surface power spectral density is, rather, maximum on a ring gago ; bobek . Stabilization of this pattern occurs when its amplitude is large enough that non-linear effects are no longer negligible. Close to the instability threshold, the rate of erosion is much smaller than the rate of addition to the surface. Hence, parameter $`ϵ`$ above, which is the ratio between these two typical rates, is small. We thus can perform a multiple scale expansion by introducing time scales $`T_1=ϵt`$ and $`T_2=ϵ^2t`$, and by rescaling length scales as $`X=ϵ^{1/2}x`$. To lowest non-linear order $`𝒪(ϵ)`$ and as seen in the slow variables ($`X`$, $`T=T_1+T_2`$), surface dynamics is described by (see tbp for details) nota3 $`_TH`$ $`=`$ $`\nu ^2H𝒦^4H+\lambda _1(H)^2\lambda _2^2(H)^2,`$ (5) where $`H=h_1+ϵh_2`$, and $`\nu `$ $`=`$ $`A\varphi ,𝒦=ϵ\gamma _0^1(D+\varphi A)[R_{eq}\gamma _0\gamma _2A(1\varphi )],`$ $`\lambda _1`$ $`=`$ $`\varphi \alpha _0(1/2\mu _2),`$ (6) $`\lambda _2`$ $`=`$ $`ϵ\alpha _0(1/2\mu _2)\gamma _0^1[(D+\varphi A)(1\varphi )R_{eq}\varphi \gamma _0\gamma _2].`$ Eq. (5) with a noise term, has been already employed in the growth of amorphous thin films raible . In our context, Eq. (5) has some important limits. First, in the absence of ion bombardment, $`A=\alpha _0=0`$, $`\gamma _0\tau ^1`$ and $`\gamma _2\mathrm{\Lambda }`$, and in the original variables (5) reduces to Mullins’ equation for thermal surface diffusion Mullins , $`_th=DR_{eq}\mathrm{\Lambda }^4h`$. In the general case, (5)-(6) include contributions to surface diffusion that are both thermally activated, and directly induced by the ion beam as in viscousflow . Second, the BH limit corresponds to $`\varphi =1`$, i.e., no redeposition. While in aste the BH limit zeroes out the $`k^4`$ contribution to the analog of $`\omega _k`$ —thus making the typical length scale of the dot structures remain undefined within linear instability—, here Eq. (5) recovers for $`\varphi =1`$ the equation obtained within BH’s approach to Sigmund’s theory cuerno\_makeev ; beyond\_ks , including the fact that the coefficients of the two nonlinear terms have the same signs thus making the equation nonlinearly unstable and mathematically ill-posed beyond\_ks ; comment . Thus, beyond its experimental relevance, redeposition is crucial in order to make the theory mathematically sound. On the other hand, the BH limit allows us to extract the phenomenological dependence of the parameters in our model with characteristics of the collision cascades, such as the ion penetration depth, $`a`$, and the longitudinal and lateral widths $`\sigma `$, and $`\mu `$, characterizing the Gaussian decay of enery deposition sigmund . Thus, for $`\varphi 1`$ we have, in the notation of cuerno\_makeev , $`\alpha _0=F`$, $`\alpha _2=a\mu ^2/(2\sigma ^2)`$, $`\mu _2=1\mu ^2/(2\sigma ^2)\mu ^2/(2\sigma ^4)`$, $`R_{eq}\gamma _2=\mu ^2/4`$, with $`FJE/\sigma `$, where $`J`$ and $`E`$ are the average ion flux and energy, respectively. Eq. (5) describes the evolution of the erosion process. Initially, dynamics is controlled by the linear terms, with the same dispersion relation $`\omega _k`$ as above, and a periodic pattern develops, with characteristic wavelength given by $$l_c=2\pi \left[2R_{eq}\gamma _2(D+\varphi A)[1ϵ(1\varphi )]/(A\varphi )\right]^{1/2},$$ (7) providing the typical size of the nanostructures that form. When local slopes become large, the nonlinear terms in Eq. (5) control the dynamics in an opposing way. While the $`\lambda _2`$ term tends to coarsen the nanostructures in amplitude and lateral size, similarly to its rôle in the coarsening of ripples on aeolian sand dunes sand , the nonlinearity $`\lambda _1`$ tends to disorder the pattern leading to the paradigmatic KS spatiotemporal chaos. Remarkably, $`\lambda _1(h)^2`$ seems to interrupt the coarsening process induced by $`\lambda _2^2(h)^2`$ and the stationary state morphology consists of domains of hexagonally ordered nanostructures separated by defects. The density of these is a function of the ratio $`r=\lambda _2/\lambda _1`$, whose $`r0`$ limit in Eq. (5) leaves us with the KS equation. In Fig. 1$`(a)`$, we plot the stationary-state morphology obtained by numerical integration of Eq. (5) for a relatively large ratio $`r=5`$ new\_note2 . The high degree of in-plane short range hexagonal ordering is made clear by the height autocorrelation function, shown in the inset of Fig. 1$`(a)`$. The time evolution of the dot pattern can be assessed in Fig. 2$`(a)`$, in which the surface roughness (mean height square deviation) $`W(t)`$ vs $`t`$ is shown for the same parameters as in Fig. 1$`(a)`$. In excellent agreement with measurements for nanodots on GaSb bobek , the roughness first increases exponentially during development of the linear instability, attains a maximum value after dots have coarsened to form a densely packed array, and finally relaxes to a smaller stationary value when defects among different dot domains are annihilated. Times between linear instability and maximum in the roughness correspond to non-linear coarsening of the dot structures, as seen in the plot of the lateral correlation length $`\xi _c(t)`$, shown on the same panel. We define $`\xi _c(t)`$ as the length-scale provided by the first secondary maximum of the height autocorrelation. As seen in Fig. 2$`(a)`$, $`\xi _c(t)`$scale is constant during linear instability, grows as $`t^{0.27\pm 0.02}`$, and saturates at long times, in agreement with experiments on InP frost . This interrupted coarsening process has been also observed on Si gago $`(b)`$ and GaSb bobek . Experimental conditions reflect in the value of $`r`$ new\_note3 , and can be such that this parameter is substantially smaller. Dynamics is then closer to that of the KS equation. The intermediate coarsening regime narrows, and is followed by kinetic roughening. A surface morphology produced in these conditions \[$`r=0.5`$\] is shown in Fig. 1$`(b)`$, which can be compared with an AFM scan \[$`(d)`$\] of a Si target irradiated as in gago . Again, agreement is excellent. Note that the morphology now differs appreciably from that of the KS equation, displayed in Fig. 1$`(c)`$. While for Eq. (5) a short-range ordered pattern coexists with long-range disorder and roughening, in the pure KS system disorder of the cellular structure is paradigmatic, see the height autocorrelations in Figs. 1$`(b),(c)`$. Still, the time evolution of the roughness in Fig. 2$`(b)`$ ($``$), predicted by Eq. (5) for small $`r`$ values, is similar to that of the KS case, Fig. 2$`(b)`$ ($`+`$): initial rapid growth is followed by much slower dynamics, and saturation to the stationary state. Such is also the experimental behavior found for nanostructures produced on Si, see Fig. 3 in gago $`(a)`$. Comparing the two plots in Fig. 2$`(b)`$, for small (non-zero) $`r`$ values the small-scale nonlinearity $`\lambda _2`$ is seen to stabilize the linear instability earlier, and leads to smaller stationary roughness. Moreover, in contrast with Fig. 2$`(a)`$, Fig. 2$`(b)`$ shows that for small or zero $`r`$ values, the roughness does not have a local maximum as a function of time. In summary, we have introduced a continuum model for the formation of nanometric sized patterns by IBS. The model accounts within an unified framework for experimental features of nanopatterns recently produced on diverse materials. Moreover, it leads to an effective interface equation providing new predictions. Thus, considering dependencies sigmund on ion energy $`E`$ of the features of the distribution of deposited energy, $`a`$, $`\mu `$, $`\sigma `$, the dot size $`l_c`$ behaves, in the large redeposition limit $`\varphi 1`$, as $`l_c[E+\mathrm{const}.]^{1/2}`$. For small $`E`$, this implies $`l_c`$ is energy independent, while $`l_cE^{1/2}`$ for large enough energies. Observations exist bobek ; frost2 compatible with such energy dependence, although a systematic study assessing the importance of redeposition would be highly desirable. From a fundamental point of view, Eq. (5) also leads to new results. Specifically, this is a height equation with local interactions in which a pattern is stabilized with constant wavelength and amplitude, in contrast with conjectures for 1$`d`$ systems coarsening . Although more theoretical work is still needed \[e.g., regarding the asymptotic properties of Eq. (5)\] this suggests that in 2$`d`$ patterns, coarsening dynamics is indeed more complex than in 1$`d`$ interrupted . ###### Acknowledgements. R. G. acknowledges a Ramón y Cajal Fellowship from MECD (Spain). This work has been partially supported by MECD (Spain) grants Nos. BFM2003-07749-C05, -05 (M. C.), -01 (R. C.), and -02 (L. V.).
warning/0506/cond-mat0506790.html	ar5iv	text	# Current-carrying molecules: a real space picture ## Abstract An approach is presented to calculate characteristic current vs voltage curves for isolated molecules without explicit description of leads. The Hamiltonian for current-carrying molecules is defined by making resort to Lagrange multipliers, while the potential drop needed to sustain the current is calculated from the dissipated electrical work. Continuity constraints for steady-state DC current result in non-linear potential profiles across the molecule leading, in the adopted real-space picture, to a suggestive analogy between the molecule and an electrical circuit. Experiments on single-molecule junctions are challenging, mainly due to the need of contacting a microscopic object, the molecule, with macroscopic leads Reed et al. (1997); Cui et al. (2001); Reichert et al. (2002); Xu and Tao (2003). Theoretical modeling of molecular junctions is difficult Nitzan and Ratner (2003); Datta (2004); Berman and Mukamel (2004); Sai et al. (2005), and again the description of contacts represents a delicate problem. To attack the complex problem of conduction through a molecular junction a strategy is emerging Kosov (2004); Burke et al. (2005) that focuses attention on isolated molecules and describes the intrinsic molecular conductivity in the absence of electrodes. At variance with common approaches that impose a potential bias to the electrodes and then calculate the resulting current Nitzan and Ratner (2003); Datta (2004); Berman and Mukamel (2004); Sai et al. (2005), a steady-state DC current is forced in the isolated molecule by making resort to a Lagrange-multiplier technique Kosov (2004), or by drawing a magnetic flux through the molecule Burke et al. (2005). Whereas the strategy is promising, two main problems remain to be solved: (1) the calculation of the potential drop needed to sustain the current, and (2) the definition of the potential profile in the molecule. Here I demonstrate that the Joule law can be used to calculate the potential drop from the electrical power dissipated on the molecule. Moreover, continuity constraints for steady-state DC current are implemented in polyatomic molecules in terms of multiple Lagrange multipliers that yield to non-linear potential profiles in the molecule. Finally, in the adopted real-space picture, the current flows through chemical bonds rather than through energy levels, leading to a suggestive description of the molecule as an electrical circuit with resistances associated to chemical bonds. To start with consider a diatomic Hubbard molecule, whose Hamiltonian $`H_0`$ is defined by $`U`$, $`t`$, and the difference of on-site energies: $`2\mathrm{\Delta }=ϵ_2ϵ_1`$. Following Kosov Kosov (2004), a current is forced through the molecule by introducing a Lagrange multiplier, $`\lambda `$, as follows: $$H(\lambda )=H_0\lambda \widehat{j}$$ (1) where $`c_{i,\sigma }^{}`$ creates an electron with spin $`\sigma `$ on the $`i`$-site, and $`\widehat{j}=it_\sigma (c_{1\sigma }^{}c_{2\sigma }H.c.)`$ measures the current flowing through the bond. Here and in the following $`\mathrm{}`$ and the electronic charge are set to 1, and $`t`$ is taken as the energy unit. The ground state of $`H(\lambda )`$, $`\|G(\lambda )`$, carries a finite current, $`J=G(\lambda )\|\widehat{j}\|G(\lambda )`$, and the Lagrange multiplier, $`\lambda `$, is fixed by imposing a predefined $`J`$ Kosov (2004). Other molecular properties can be calculated as well, and their dependence on $`J`$ can be investigated Kosov (2004). Just as an example, the bond-order decreases with $`J`$, and, in systems with inequivalent sites, the on-site charge distribution is equalized by the current flow. These are interesting informations, but characteristic $`J(V)`$ curves are still needed. The Lagrange multiplier, $`\lambda `$, has the dimensions and the meaning of a magnetic flux drawn across the molecule to generate a spatially uniform electric field Kohn (1964); Burke et al. (2005), $`E\omega \lambda `$, where $`\omega `$ is the field frequency Kohn (1964). In the limit of static fields, $`\omega 0`$, both $`E`$ and $`V`$ vanish, suggesting that a finite current flows in the molecule at zero bias. This contrasts sharply with the fundamental relation between charge transport and energy dissipation Nitzan (2001); Datta (2004); Burke et al. (2005): a finite $`V`$ is needed to sustain a current due to dissipative phenomena occurring in the conductor. Specifically, the Joule law relates the potential drop in a conductor to the electrical power spent on the system to sustain the current, $`W=VJ`$. Since $`J`$ is known, $`V`$ can be obtained from a calculation of the dissipated power. Dissipation is conveniently described in the density matrix formalism using as a basis of the eigenstates $`\|k`$ of $`H(0)`$ Mukamel (1995); Boyd (2003). The equilibrium density matrix, $`\sigma _0`$, is a diagonal matrix whose elements are fixed by the Boltzmann distribution. On the same basis $`\sigma (\lambda )`$ is a non-diagonal matrix corresponding to a non-equilibrium state whose dynamics is governed by: $`\dot{\sigma }=\frac{i}{\mathrm{}}[H,\sigma ]+\dot{\sigma }_R`$, where $`\dot{\sigma }_R`$ accounts for relaxation phenomena, as due to all degrees of freedom not explicitly described by $`H`$ (e.g. molecular vibrations, or environmental degrees of freedom also including leads) Mukamel (1995); Datta (2004); Nitzan (2001). Diagonal elements of $`\dot{\sigma }_R`$ describe depopulation and are associated with energy dissipation. As for depopulation I adopt a simple phenomenological model with $`(\dot{\sigma }_R)_{kk}=_m\gamma _{km}\sigma _{mm}`$, where $`\gamma _{km}`$ measures the probability of the transition from $`k`$ to $`m`$ Mukamel (1995); Boyd (2003). For the sake of simplicity I will consider the low-temperature limit, with $`\gamma _{km}=\gamma `$ for $`k>m`$ and $`\gamma _{km}=0`$ otherwise. The dynamics of off-diagonal elements is governed by depopulation and dephasing effects: $`(\dot{\sigma }_R)_{km}=\mathrm{\Gamma }_{km}\sigma _{km}`$, with $`\mathrm{\Gamma }_{km}=(\gamma _{kk}+\gamma _{mm})/2+\gamma _{km}^{}`$, where $`\gamma _{kk}=_m^{}\gamma _{mk}`$, and $`\gamma _{km}^{}`$ describes dephasing, i.e. the loss of coherence due to elastic scattering Mukamel (1995); Boyd (2003). The energy dissipated by the system, $`Tr(\dot{\sigma }_rH)`$, has two contributions: the first one, $`W_d=_k(\dot{\sigma }_R)_{kk}`$, is always negative and measures the energy that the system dissipates to the bath as the current flows. This term is governed by depopulation, whereas dephasing plays no role. The second term, $`W=\lambda Tr(\dot{\sigma }_R\widehat{j})`$, measures the electric work done on the system to sustain the current: it is this term that enters the Joule law. In non-degenerate systems $`\widehat{j}`$ is an off-diagonal operator, so that only off-diagonal elements of $`\dot{\sigma }_R`$ enter the expression for $`W`$. Both depopulation and dephasing then contribute to $`V`$, and hence to the molecular resistance. This is in line with the observation that a current flowing through a molecule implies an organized motion of electrons along a specific direction Buttiker (1985); Nitzan (2001); Datta (2004). Therefore any mechanism of scattering, either anelastic, as described by depopulation, or elastic, as described by dephasing, contributes to the electrical resistance Buttiker (1985); Datta (2004). The unbalance between $`W_d`$ and $`W`$ is always positive: the molecule heats as current flows. Efficient heat dissipation is fundamental to reach a steady-state regime and to avoid molecular decomposition Nitzan (2001). Fig. 1 shows the characteristic curves calculated for a diatomic molecule with $`\gamma =0.2`$ and $`\gamma _{km}^{}=0`$. In the left panel results are shown for the symmetric, $`\mathrm{\Delta }=0`$, system. As expected, electronic correlations decrease the conductivity. The results in the right panel for an asymmetric system ($`\mathrm{\Delta }0`$) show instead an increase of the low-voltage conductivity with increasing $`U`$. This interesting result is related to the minimum excitation gap, and hence the maximum conductance, of the system with $`U=2\mathrm{\Delta }`$. The asymmetric diatomic molecule represents a minimal model for the Aviram and Ratner rectifier Aviram and Ratner (1974), however the characteristic curves in the right panel of Fig. 1 are symmetric, and do not support rectification. In agreement with recent results, rectification in asymmetric molecules is most probably due to contacts Datta (2004), or to the coupling between electrons and vibrational or conformational degrees degrees of freedom Troisi and Ratner (2002). Before attacking the more complex problem of polyatomic molecules, it is important to compare the results obtained so far with well known results for the optical conductivity Kohn (1964). If $`\mathrm{\Gamma }_{km}=\mathrm{\Gamma }`$, as it occurs, e.g., for systems with large inhomogeneous broadening (the coherent conductance limit Nitzan and Ratner (2003); Nitzan (2001)), the expression for the potential drop is very simple: $`V=\lambda \mathrm{\Gamma }`$. Then, a perturbative expansion of $`J`$ leads to the following expression for the zero-bias conductivity, $`𝒢_0`$: $$𝒢_0=\frac{2}{\mathrm{\Gamma }}\underset{k}{}\frac{k\|\widehat{j}\|g\|^2}{E_kE_g}$$ (2) where $`g`$ is the gs of $`H(0)`$, with energy $`E_g`$, and the sum runs on all excited states. This expression for the DC conductivity coincides with the zero-frequency limit of the optical conductivity Kohn (1964), provided that the frequency, $`\omega `$, appearing in the denominator of the expression for the optical conductivity in Ref. Kohn (1964) is substituted by $`\omega i\mathrm{\Gamma }`$. Introducing a complex frequency to account for relaxation is a standard procedure in spectroscopy Boyd (2003), leading to similar effects as the introduction of an exponential switching on of the electromagnetic field Kohn (1964): both phenomena account for the loss of coherence of electrons driven by an EM field and properly suppress the divergence of the optical conductivity due to the build-up of the phase of electrons driven by a static field. The connection between DC and optical conductivity breaks down in polyatomic molecules. To keep the discussion simple, I will focus attention on linear Hubbard chains. The optical conductivity of Hubbard chains was discussed based on the current operator $`\widehat{J}=ie_it_i(c_{i,\sigma }^{}c_{i+1,\sigma }H.c.)/\mathrm{}`$ Maldague (1977). However, this operator measures the average total current and does not apply to DC currents. Specifically, if a term $`\lambda \widehat{J}`$ is added to the molecular Hamiltonian, a finite average current is forced through the molecule Kosov (2004), but this current does not satisfy basic continuity constraints for steady-state DC current. In fact, to sustain a steady-state DC current one must avoid the build up of electrical charge at atomic sites. Specifically, in linear molecules the continuity constraint imposes that exactly the same amount of current flows through each bond in the molecule. To impose this constraint the current on each single bond must be under control and a Lagrange multiplier must be introduced for each bond, as follows: $$H(\lambda _i)=H_0\underset{i}{}\lambda _i\widehat{j}_i$$ (3) where $`\widehat{j}_i=it_i(c_{i,\sigma }^{}c_{i+1,\sigma }H.c.)/\mathrm{}`$, and the $`\lambda _i`$’s are fixed by imposing $`j_i=G\|\widehat{j}_i\|G=J`$ independent on $`i`$. As before, the electrical work done on the molecule is: $$W=\underset{i}{}\lambda _iTr(\widehat{j}_i\dot{\sigma }_R)$$ (4) that naturally separates into contributions, $`W_i`$, relevant to each bond. The total potential drop across the molecule, $`V=W/J`$, is then the sum of the potential drops across each bond, $`V_i=W_i/J`$, leading in general to non-linear potential profiles. Of course, in the adopted real-space picture the potential profile can only be calculated at atomic positions, and no information can be obtained on the potential profile inside each bond. Therefore, instead of showing the potential profile along the molecule, I prefer to convey the same information in terms of bond-resistances, defined as: $`R_i=(J/V_i)^1`$. Fig. 2 shows the behavior of a 3-site chain with three electrons, $`U=4`$, equal on-site energies and different $`t`$: $`t_1=1.2`$, and $`t_2=0.8`$. The left panel shows the characteristic $`J(V)`$ curve, and continuous lines in the right panel report the total resistance $`R`$, and the two bond resistances, $`R_1`$ and $`R_2`$. The molecular resistance varies with the applied voltage and, as expected, the resistance of the weaker bond is higher than the resistance of the stronger bond. Dimensionless resistances in the figure are in units with $`\mathrm{}/e^2=(2\pi \stackrel{~}{g}_0)^1=1`$, where $`\stackrel{~}{g}_0`$ is the quantum of conductance, that, in standard approaches to molecular junctions Datta (2004) represents the maximum conductance associated with a discrete molecular level. This well known result is related to the inhomogeneous broadening of molecular energy levels as due to their interaction with the electrodes Datta (2004). Of course there is no intrinsic limit to the conductivity in the model for isolated molecules discussed here. As a direct consequence of the continuity constraint, the total resistance $`R=(J/V)^1`$ is the sum of the two bond-resistances, leading to a suggestive description of the molecule as an electrical circuit, with resistances associated with chemical bonds joint in series at atomic sites. Whereas this picture is useful, the concept of bond-resistance should be considered with care in molecular circuits. At variance with standard conductors, in fact, the resistance of the bonds depends not only on the circuit (the molecule) they are inserted in, but also on the way the resistance is measured. Dashed lines in Fig. 2 show the bond-resistances calculated by forcing the current through specific bonds (i.e. by setting a single $`\lambda _i0`$ in Eq. 3), and these differ from the bond-resistances calculated when the whole molecule carries the current (continuous lines). The situation becomes somewhat simpler in the coherent conductance limit, $`\mathrm{\Gamma }_{km}=\mathrm{\Gamma }`$. Perturbative arguments can be used to demonstrate that at zero bias the bond resistances calculated for the current flowing through the whole molecule or through a single bond do coincide. This additive result for the molecular resistance in the coherent transport limit is in line with the observation of transmission rates inversely proportional to the molecular length in the same limit Davis et al. (1997). However, as shown in Fig. 3, this simple Ohmic behavior breaks down quickly at finite bias. The introduction of as many Lagrange multipliers as many current-channels (bonds) are present in the molecule accounts for a non-linear potential profile through the molecule, i.e. for a non-uniform electric field. Accounting for a single Lagrange multiplier coupled to the total current operator is equivalent to draw a magnetic flux through the molecule as to generate a spatially homogeneous electric field Kohn (1964); Maldague (1977); Burke et al. (2005), a poor approximation for DC conductivity in extended (polyatomic) molecules. Just as an example, for a 4-site chain with the same $`t_i=1`$ on each bond, the zero-bias resistance of the central bond exceeds that of the lateral bonds with $`R_2/R_1`$ ranging from 10 to 2 as $`U`$ increases from 0 to 4. Bonds with the same $`t`$ have different resistances due to their different bond-orders, and, in agreement with recent results Liang et al. (2004); Berman and Mukamel (2004), this demonstrates nicely the need of accounting for non-uniform electric fields in extended molecules, even for very idealized molecular structures. It is of course possible to discuss more complex molecular models. As an interesting example, a nearest-neighbor hopping $`t^{}`$ is added to the Hamiltonian for the three site molecule discussed above. This opens a new channel for electrical transport, and a term $`\lambda ^{}\widehat{j}^{}`$ adds to the Hamiltonian with $`\widehat{j}^{}=it^{}_\sigma (c_{1\sigma }^{}c_{3\sigma }H.c.)`$. As before, continuity imposes $`j_1=j_2`$, as to avoid building up of charge at the central site. The total current is $`J=j_1+j^{}`$ and, of course, no continuity constraint is given on $`j^{}`$. However the potential drop across the molecule, i.e. the potential drop measured at sites 1 and 3 must be uniquely defined. Therefore one must tune $`\lambda _1`$, $`\lambda _2`$ and $`\lambda ^{}`$ as to satisfy $`j_1=j_2`$, while satisfying the condition: $`V_1+V_2=V^{}=V`$, with $`V_i=W_i/j_i`$ and $`V^{}=W^{}/j^{}`$. Imposing a constraint on the potentials is a tricky affair, that becomes trivial when $`\mathrm{\Gamma }_{km}=\mathrm{\Gamma }`$. In that case in fact $`V_i=\mathrm{\Gamma }\lambda _i`$ and $`V^{}=\mathrm{\Gamma }\lambda ^{}`$, so that the constraint on the potentials immediately translates into a constraint on Lagrange multipliers. Fig. 4 shows some results obtained in this limit for a system with $`t_1=t_2=1`$, $`t^{}=0.4`$. In spite of the fairly large $`t^{}`$ value, the contribution to the current from the bridge-channel is small, mainly due to the small bond-order for next-nearest neighbor sites. Once again the physical constraints imposed to the currents and to the potentials lead to standard combination rules for bond-resistances with $`1/R=1/R^{}+1/(R_1+R_2)`$. As for the DC conductivity is concerned, the molecule behaves as an electrical circuit with two resistances, $`R_1`$ and $`R_2`$ in series bridged by a parallel resistance, $`R^{}`$. Applying the proposed approach to complex molecular structures and/or to molecules described by accurate quantum chemical Hamiltonians is non-trivial due to the appearance in the Hamiltonian of as many Lagrange multipliers as many current channels are considered, and due to the large number of constraints to be implemented. Instead, at least for small molecules, the approach can be fairly easily extended to account for vibrational degrees of freedom. Non-adiabatic calculations are currently in progress to describe electrical conduction through a diatomic molecule in the presence of Holstein and Peierls electron-phonon coupling. More refined models for the relaxation dynamics can also be implemented Davis et al. (1997), whereas the introduction of spin-orbit coupling can lead to a model for spintronics. In conclusion, this paper presents an approach to the calculation of characteristic current/voltage curves for isolated molecules in the absence of contacts. While hindering the direct comparison with experimental data, this allows the definition of the molecular conductivity as an intrinsic molecular property. Even more important, a paradigm is defined for imposing a steady-state DC current through the molecule, while extracting the voltage drop across the molecule from the energy dissipation. The careful implementation of continuity constraints for steady-state DC current leads to the definition of the potential profile through the molecule, that, in the adopted real-space description, quite naturally results in the concept of bond-resistances, in a suggestive description of the molecule as an electrical circuit with current flowing through chemical bonds. I thank D. Kosov for useful discussions and correspondence. Discussions with A. Girlando, S. Pati, S. Ramasesha, and Z.G. Soos are gratefully acknowledged. The contributions from S. Cavalca and C. Sissa in the early stages of this work are acknowledged. Work supported by Italian MIUR through FIRB-RBNE01P4JF and PRIN2004033197-002.
warning/0506/astro-ph0506102.html	ar5iv	text	# The radial velocity dispersion profile of the Galactic halo: Constraining the density profile of the dark halo of the Milky Way ## 1 Introduction The determination of the total mass of the Galaxy has been a subject of considerable interest since the work of Kapteyn in the early 1920s (see Fich & Tremaine 1991 for a nice introductory review on the subject). Since then, the mass of the Milky Way has seen its estimates grow by factors of ten to a hundred, with some dependence on the type of mass tracer used: H I kinematics, satellite galaxies and globular clusters, or the Local Group infall pattern. The most recent determinations yield fairly consistent values for the mass within 50 kpc, with an uncertainty of the order of 20% for a given mass model (Kochanek 1996; Wilkinson & Evans 1999, hereafter W&E99; Sakamoto, Chiba & Beers 2003, hereafter SCB03). However, even today, the total mass of the Galaxy is not known better than within a factor of two. Whatever method is used, be it the H I kinematics, globular clusters, satellite galaxies, or halo giants, it is only possible to determine the mass enclosed in the region probed by these tracers (Binney & Tremaine 1987). This implies that the rotation curve derived from H I will only constrain the mass within roughly 18 kpc from the Galactic centre (Rohlfs & Kreitschmann 1988; Honma & Sofue 1997), a region which is baryon dominated. Globular clusters and satellite galaxies are, in principle, better probes of the large scale mass distribution of the Galaxy, since they are found out to distances beyond 100 kpc. However, there are only 15 such objects beyond 50 kpc (Zaritsky et al. 1989; Kochanek 1996). Only 6 of these have proper motion measurements, which despite the large errors, can further constrain the shape of the velocity ellipsoid. Using this dataset, W&E99 favour isotropic to slightly tangentially anisotropic models, although 1$`\sigma `$ contours for the velocity anisotropy $`\beta `$ give $`0.4\beta 0.7`$. SCB03 have added to the sample used by W&E, field blue horizontal branch stars with proper motions and radial velocities. While this is clearly an improvement, these stars are located within 10 kpc of the Sun, which strongly limits their constraining power at larger radii. In their models, the velocity ellipsoid is tangentially anisotropic, with $`\beta 1.25`$ as the most likely value. It is clearly important to measure the total mass of the Galaxy in order to constrain its dark-matter content. However, it is also critical to determine its distribution: density profile, flattening, velocity ellipsoid, etc. One of the most fundamental predictions of cold-dark matter models is that the density should follow an NFW profile throughout most of the halo (Navarro, Frenk & White 1997). The density profiles derived from the gas rotation curves of large samples of external galaxies do not always follow the NFW shape (de Blok et al. 2001). Tracers at larger distances are rare, but objects such as planetary nebulae or globular clusters could yield powerful constrains on the mass distribution at those radii, for example for elliptical galaxies as shown by Romanowsky et al. (2003). In the case of the Milky Way, the situation is not dissimilar. The distribution of mass inside the Solar circle has been studied extensively (see e.g. Dehnen & Binney 1998; Evans & Binney 2001; Bissantz, Debattista & Gerhard 2004). A common conclusion is that there is little room for dark-matter in this region of the Galaxy. But does the dark-matter beyond the edge of the Galactic disk follow an NFW profile? How does the most often assumed isothermal profile perform in this region of the Galaxy (e.g. Sommer-Larsen et al. 1997; Bellazzini 2004)? Is the velocity ellipsoid close to isotropic as found in CDM simulations (Ghigna et al. 1998)? Modeling of the kinematics of halo stars by Sommer-Larsen et al. (1997) favoured an ellipsoid that became more tangentially anisotropic towards larger distances, while Ratnatunga & Freeman (1989) found a constant line-of-sight velocity dispersion out to 25 kpc. These fundamental issues can only be addressed when a sufficiently large number of probes of the outer halo of the Galaxy are available. Ideal tracers are red giant stars or blue horizontal branch stars, which can be identified photometrically also at large galactocentric distances (Morrison et al. 2000; Clewley et al. 2002; Sirko et al. 2004a,b). Spectroscopic follow-up allows both the confirmation of the luminosity class as well as the determination of radial velocities with relatively small errors (Morrison et al. 2003). With the advent of wide field surveys, such as the Sloan Digital Sky Survey, or the Spaghetti survey, the numbers of such outer halo probes have increased by large amounts, making this an ideal time to address the mass distribution of our Galaxy in greater detail. This paper is organized as follows. In the next section we describe the observational datasets used to determine the radial velocity dispersion curve. In Sec. 2.2 we introduce several mass models for the dark halo of our Galaxy and derive how the line of sight velocity dispersion depends on the model parameters. In Sec. 2.3 we compare the data to the models and derive the best fit values of the parameters using $`\chi ^2`$ fitting. Finally we discuss our results and future prospects in Sec. 3. ## 2 The radial velocity dispersion curve ### 2.1 The observational datasets Our goal is to derive the radial velocity dispersion profile of the Milky Way stellar halo in the regime where it is dominated by the gravitational potential of its dark-matter halo. Hence we restrict ourselves to tracers located at Galactocentric distances greater than 10 kpc, where the disc’s contribution is less important. We use a sample of 9 satellite galaxies, 44 globular clusters, 57 halo giants and 130 field blue horizontal branch stars (FHB). The various data sources of this sample are listed in Table 1. It is worth noting that there are 24 objects located beyond 50 kpc in our sample, and that we have enough statistics to measure radial velocity dispersion out to 120 kpc as shown in the top panel of Figure 1. This covers a significantly larger radial range than many previous works, including e.g. Sommer-Larsen et al. (1997), whose outermost point is at 50 kpc. The red halo giants are from the “Spaghetti” Survey (Morrison et al. 2000). This is a pencil beam survey that has so far covered 20 $`\mathrm{deg}^2`$ in the sky, down to $`V20`$. It identifies candidate halo giants using Washington photometry, where the 51 filter<sup>1</sup><sup>1</sup>1The 51 filter is centered on the Mgb/MgH feature near 5170 $`\dot{\mathrm{A}}`$ allows for a first luminosity selection. Spectroscopic observations are then carried out to confirm the photometric identification and to determine the radial velocities of the stars. We have derived the heliocentric distance for the FHB stars from Wilhelm et al. (1999b) using the relation $$M_V(\mathrm{HB})=0.63+0.18([\mathrm{Fe}/\mathrm{H}]+1.5)$$ (Carretta et al. 2000). In all cases, accurate distances and radial velocities are available: the average error in velocity ranges from a few km s<sup>-1</sup> (satellite galaxies and globular clusters) to 10-15 km s<sup>-1</sup> (FHB stars and red giants); the typical relative distance error is approximately 10%. When transforming the heliocentric l.o.s. velocities, $`V_{\mathrm{los}}`$, into Galactocentric ones, $`V_{\mathrm{GSR}}`$, we assume a circular velocity of $`V_{\mathrm{LSR}}=220`$ km s<sup>-1</sup> at the solar radius ($`R_{}=`$ 8 kpc) and a solar motion of ($`U`$,$`V`$,$`W`$) $`=`$ (10, 5.25, 7.17) km s<sup>-1</sup> , where $`U`$ is radially inward, $`V`$ positive in the direction of the Galactic rotation and $`W`$ towards the North Galactic Pole (Dehnen & Binney 1998). Hereafter we refer to: the radial velocity (dispersion) measured in a heliocentric coordinate system as the l.o.s. velocity, $`V_{\mathrm{los}}`$ (dispersion, $`\sigma _{\mathrm{los}}`$); the l.o.s. velocity (and its dispersion) corrected for the solar motion and the LSR motion as the Galactocentric radial velocity, $`V_{\mathrm{GSR}}`$ (dispersion, $`\sigma _{\mathrm{GSR}}`$); the radial velocity (and its dispersion) in a reference frame centered on the Galactic Centre as the true radial velocity, $`V_r`$ (dispersion, $`\sigma _r`$). Figure 2 shows $`V_{\mathrm{GSR}}`$ as function of the Galactocentric distance $`r`$ for all the objects used in this work. The bottom panel in Figure 1 shows the Galactocentric radial velocity dispersion as function of distance from the Galactic centre. This is computed in bins whose width is approximately twice the average distance error of objects in the bin. This implies that our bin sizes range from 3 kpc at $`r10`$ kpc, to 40 kpc at $`r120`$ kpc. The error-bar on the velocity dispersion in each bin is calculated performing Monte Carlo simulations. We assume the velocity and distance errors are gaussianly distributed in the heliocentric reference frame. In practice, this means that we randomly generate velocities and distances for each one of the stars, whose mean and dispersion are given by the observed value and its estimated error, respectively. We then convert the heliocentric quantities into Galactocentric ones. We repeat this exercise for 10,000 sets, and for each of these we measure $`\sigma _{\mathrm{GSR}}`$ in the same bins as the original data. We use the rms of this velocity dispersion, obtained from the 10,000 simulations, as the error on the velocity dispersion we measured in the bin. One may question whether the satellite galaxies can be considered fair tracers of the gravitational potential of the dark matter halo of the Milky Way (e.g. Taylor, Babul & Silk 2004; Gao et al. 2004). To get a handle on this issue, we compute the velocity dispersion profile both with and without them (squares and diamonds, respectively in Fig. 1). Since the trend is similar in both cases we may consider the satellites to be reliable probes of the outer halo potential. ### 2.2 The models #### 2.2.1 Jeans equations If we assume that the Galactic halo is stationary and spherically symmetric we can derive the (expected) radial velocity dispersion profile $`\sigma _{r,}`$ of the stars from the Jeans equation (Binney & Tremaine 1987): $$\frac{1}{\rho _{}}\frac{d(\rho _{}\sigma _{r,}^2)}{dr}+\frac{2\beta \sigma _{r,}^2}{r}=\frac{d\varphi }{dr}=\frac{V_\mathrm{c}^2}{r}$$ (1) where $`\rho _{}(r)`$ is the mass density of the stellar halo, $`\varphi (r)`$ and $`V_\mathrm{c}(r)`$ are the potential and circular velocity of the dark matter halo and $`\beta `$ is the velocity anisotropy parameter, defined as $`\beta =1{\displaystyle \frac{\sigma _\theta ^2}{\sigma _r^2}}`$, and assuming $`\sigma _\theta ^2=\sigma _\varphi ^2`$. Note that $`\beta =0`$ if the velocity ellipsoid is isotropic, $`\beta =1`$ if the ellipsoid is completely aligned with the radial direction, while $`\beta <0`$ for tangentially anisotropic ellipsoids. The Jeans equation allows us to determine a unique solution for the mass profile if we know $`\sigma _{r,}^2(r)`$, $`\rho _{}(r)`$ and $`\beta (r)`$, although this solution is not guaranteed to produce a phase-space distribution function that is positive everywhere. We are, however, faced with two uncertainties: the velocity anisotropy and the behaviour of the stellar halo density at very large distances. The latter has been determined to vary as a power-law $`\rho _{}(r)r^\gamma `$ with $`\gamma `$ 3.5 out to $`50`$ kpc (Morrison et al. 2000; Yanny et al. 2000), and we shall assume this behaviour can be extrapolated all the way out to our last measured point. More crucial is the unknown variation of the velocity anisotropy with radius, which is difficult to determine because of the lack of tracers with accurate proper motions beyond the Solar neighbourhood. This implies in principle, that large amounts of kinetic energy can be hidden to the observer, an effect known as the mass-velocity anisotropy degeneracy. For sake of simplicity, and given that the situation is unlikely to change until the advent of new space astrometric missions such as SIM and Gaia (Perryman et al. 2001), throughout most of this work we shall make the assumption that $`\beta `$ is constant, i.e. independent of radius $`r`$. To derive Eq. (1) we have assumed that the stellar halo can be considered as a tracer population of objects moving in an underlying potential. This is justified by the negligible amount of mass present in this component, compared to, for example, that in the disk and the dark halo. The (expected) radial velocity dispersion for the tracer population $`\sigma _{r,}`$ may be thus derived by integrating Eq. (1). This leads to $$\sigma _{r,}^2(r)=\frac{1}{\rho _{}e^{{\scriptscriptstyle 2\beta 𝑑x}}}_{x}^{}{}_{}{}^{\mathrm{}}\rho _{}V_\mathrm{c}^2e^{{\scriptscriptstyle 2\beta 𝑑x^{\prime \prime }}}𝑑x^{},x=\mathrm{ln}r.$$ (2) Here, we have used that $`r^{2\beta }\rho _{}\sigma _{r,}\|_{\mathrm{}}=0`$. Note that the radial velocity dispersion of the tracer population depends on the particular form of the circular velocity of the underlying (gravitationally dominant) mass distribution. Since proper motions are not available for the whole sample and we only have access to heliocentric velocities, the quantity that we measure is not the true radial velocity dispersion but $`\sigma _{\mathrm{GSR},}`$. When comparing this quantity to model predictions, we must take in account a correction factor for the lack of information on the tangential component of the velocity. Following the procedure described in Appendix A, we find that the Galactocentric radial velocity dispersion, $`\sigma _{\mathrm{GSR},}`$, is related to the true radial velocity dispersion, $`\sigma _{r,}`$ as $$\sigma {}_{\mathrm{GSR},}{}^{}(r)=\sigma _{r,}(r)\sqrt{1+2(1\beta )H(r)},$$ (3) where $$H(r)=\frac{r^2+R_{}^2}{4r^2}\frac{(r^2R_{}^2)^2}{8r^3R_{}}\mathrm{ln}\frac{r+R_{}}{rR_{}}.$$ (4) The above equation for $`H(r)`$ is valid at Galactocentric distances $`r>R_{}`$. For a purely radial anisotropic ellipsoid ($`\beta =`$ 1) $`\sigma _{\mathrm{GSR},}`$ and $`\sigma _{r,}`$ coincide. For a tangentially anisotropic stellar halo, the correction factor becomes negligible at distances larger than about 30-40 kpc. #### 2.2.2 Specifing dark-matter halo models We adopt three different models for the spherically symmetric dark-matter halo potential: * Pseudo-Isothermal sphere. This model has been extensively used in the context of extragalactic rotation curve work. The density profile and circular velocity associated to a pseudo-isothermal sphere are: $$\rho (r)=\rho _0\frac{r_\mathrm{c}^2}{(r_\mathrm{c}^2+r^2)},$$ (5) and $$V_\mathrm{c}^2(r)=V_\mathrm{c}^2(\mathrm{})\left(1\frac{r_\mathrm{c}}{r}\mathrm{arctg}\frac{r}{r_\mathrm{c}}\right),$$ (6) where $`r_\mathrm{c}`$ is the core radius, and $`\rho _0={\displaystyle \frac{V_\mathrm{c}^2(\mathrm{})}{4\pi Gr_\mathrm{c}^2}}`$. We set $`V_\mathrm{c}(\mathrm{})=`$ 220 km s<sup>-1</sup> as asymptotic value of the circular velocity. At large radii the density behaves as $`\rho r^2`$ giving a mass that increases linearly with radius. * NFW model. In this case the dark matter density profile is given by $$\rho (r)=\frac{\delta _c\rho _\mathrm{c}^0}{(r/r_\mathrm{s})(1+r/r_\mathrm{s})^2}$$ (7) where $`r_\mathrm{s}`$ is a scale radius, $`\rho _\mathrm{c}^0`$ the present critical density and $`\delta _\mathrm{c}`$ a characteristic overdensity. The latter is defined by $`\delta _\mathrm{c}={\displaystyle \frac{100c^3g(c)}{3}}`$, where $`c=r_\mathrm{v}/r_\mathrm{s}`$ is the concentration parameter of the halo and $`g(c)={\displaystyle \frac{1}{\mathrm{ln}(1+c)c/(1+c)}}`$. The circular velocity associated with this density distribution is $$V_\mathrm{c}^2(s)=\frac{V_\mathrm{v}^2g(c)}{s}\left[\mathrm{ln}(1+cs)\frac{cs}{1+cs}\right]$$ (8) where $`V_\mathrm{v}`$ is the circular velocity at the virial radius $`r_\mathrm{v}`$ and $`s=r/r_\mathrm{v}`$. The concentration $`c`$ has been found to correlate with the virial mass of the halo (Navarro, Frenk & White 1997; Bullock et al. 2001; Wechsler et al. 2002). However, the relation presents a large scatter. For example, for a halo of mass 1.0$`\times 10^{12}h^1M_{}`$ the predicted concentration ranges between 10 and 20. Hence, we cannot consider the NFW density profile as a one-parameter family; we need to describe it by the concentration $`c`$, and by the virial mass or the circular velocity at the virial radius. At large radii (for $`rr_\mathrm{s}`$), the density behaves as $`\rho r^3`$, and therefore, the total mass diverges logarithmically. However, we can impose that the particles must be bound at the virial radius, and so when integrating Eq. (2), we set the upper integration limit to $`r_\mathrm{v}`$ and we use $`r^{2\beta }\rho _{}\sigma _{r,}\|_{r_\mathrm{v}}=0`$. * Truncated Flat model. This density profile was recently introduced by W&E99 to describe the dark matter halo of Local Group galaxies. It is a mathematically convenient extension of the Jaffe (1983) model. The form of the density profile of the Truncated Flat model (hereafter TF) is $$\rho (r)=\frac{M}{4\pi }\frac{a^2}{r^2(r^2+a^2)^{3/2}}$$ (9) where $`a`$ is the scale length and $`M`$ the total mass of the system. For $`ra`$, the density falls off as $`\rho r^5`$. The circular velocity due to this density distribution is $$V_\mathrm{c}^2(r)=\frac{V_0^2a}{(r^2+a^2)^{1/2}}.$$ (10) We set $`V_0=`$ 220 km s<sup>-1</sup> (W&E99). The resulting rotation curve is flat in the inner part, with amplitude $`V_0=\sqrt{GM/a}`$, and becomes Keplerian for $`ra`$. Having fixed the amplitude of the circular velocity ($`V_0`$), this model is reduced to a one parameter-family characterized by the scale length $`a`$, or the mass $`M`$. ### 2.3 Results #### 2.3.1 Models with constant velocity anisotropy The methodology we use consists in comparing the measured Galactocentric radial velocity dispersion $`\sigma _{\mathrm{GSR},}`$ for each of the distance bins with that predicted for the different models discussed in Sec. 2.2. For the latter, we explore the space of parameters which define each model and determine the $`\chi ^2`$ as: $$\chi ^2=\underset{i=1}{\overset{Nbins}{}}\left(\frac{\sigma _{\mathrm{GSR}_i,}\sigma _{GSR,}(r_i;\beta ,p)}{ϵ_r}\right)^2.$$ (11) Here, the variable $`p`$ denotes a characteristic parameter of each model (e.g. scale length or total mass), while $`ϵ_r`$ is the error in the observed radial velocity dispersion as estimated through the bootstrap sampling technique described before. The best-fitting parameters are defined as those for which $`\chi ^2`$ is minimized. In the case of the isothermal sphere, the free parameters are the dark matter halo core radius, $`r_\mathrm{c}`$, and the stellar velocity dispersion anisotropy parameter, $`\beta `$. The left panel of Fig. 3 shows the $`\chi ^2`$ contours for this model. The minimum $`\chi ^2`$ value is $`\chi _{\mathrm{min}}^2=23`$ for $`r_\mathrm{c}`$ = 1.6 kpc and $`\beta =0.4`$, with 1-$`\sigma `$ contours encompassing $`0.6r_\mathrm{c}2.6`$ and $`0.7\beta 0.1`$. This corresponds to a best-fitting mass $`M=1.3\times 10^{12}`$ $`M_{}`$ (note that, since the mass for the pseudo-isothermal model is not finite, we quote the mass within our last measured point, at $`r=`$ 120 kpc). The 1-$`\sigma `$ errors on the mass, calculated from the 1-$`\sigma `$ errors for the core radius, lead to a relative error of the order of 1%. The reason for this small value is due to the fact that the best-fitting core radius is very small, and hence variations in its value (even by 100%) will barely affect the mass enclosed at large radii. On the right panel of Fig. 3 we plot the Galactocentric radial velocity dispersion for this best-fitting model. As expected, this model predicts a velocity dispersion that is roughly constant with radius. However, the observed $`\sigma _{\mathrm{GSR},}`$ shows a rather strong decline at large radii, which is not reproduced by the pseudo-isothermal halo model. The top panels of Fig. 4 show the $`\chi ^2`$ contours for the NFW model for 4 different concentrations ($`c=10`$, 14, 16, and 18). Note that the minimum $`\chi ^2`$ value decreases for increasing concentrations. Since the concentration is defined as $`c=r_\mathrm{v}/r_\mathrm{s}`$, for a fixed mass (or virial radius $`r_\mathrm{v}`$) a larger $`c`$ implies a smaller scale radius. This results in a radial velocity dispersion that starts to decline closer to the centre in comparison to a halo of lower concentration, reproducing better the trend observed in the data. Our $`\chi ^2`$ fitting technique yields for $`c=10`$ a best-fitting virial mass of 1.2$`\times 10^{12}`$ $`M_{}`$ ($`\chi _{\mathrm{min}}^2=36`$), while for $`c=18`$, $`M_\mathrm{v}=0.8\times 10^{12}`$ $`M_{}`$ ($`\chi _{\mathrm{min}}^2=12`$). We find that the velocity anisotropy for the minimum $`\chi ^2`$ is almost purely radial in all cases. In the bottom panel of Fig. 4 we show the observed Galactocentric radial velocity dispersion overlaid on two of the best-fitting NFW models. Note that beyond 40 kpc, the model with $`c=10`$ is clearly inconsistent with the data at the 1$`\sigma `$ level at $`r`$ 40 and 50 kpc and at the 2$`\sigma `$ level in the last two bins. On the other hand, the $`c=18`$ model gives a good fit of the data out to 30 kpc but overpredicts the velocity dispersion at large radii at the 1$`\sigma `$ level. We thus consider the NFW model with $`M_\mathrm{v}=0.8_{0.2}^{+1.2}\times 10^{12}`$$`M_{}`$ and $`c=18`$ as producing the best fit. Fig. 4 also shows the favourite model of Klypin et al. (2002) with $`M_\mathrm{v}=1.0\times 10^{12}`$ and $`c=`$12 (dotted curve). Since no velocity anisotropy was given in the source we performed a $`\chi ^2`$ fit to our data using the parameters from Klypin et al. (2002) and leaving $`\beta `$ as a free parameter. This favoured once again an almost purely radial anisotropy. The fit obtained in this case is very similar to that found in our $`c=10`$ model. Since our last measured point is at $`r_{\mathrm{last}}`$ 120 kpc, the constraining power of our data is stronger in the region enclosed by this radius. The value of the virial mass we just derived is an extrapolation of the model at larger distances. For completeness, we quote here the mass within 120 kpc for our best fitting NFW model with $`c=18`$, $`M(<120\mathrm{kpc})=5.4_{1.4}^{+2.0}\times 10^{11}`$$`M_{}`$ (the errors are calculated from the 1$`\sigma `$ errors in the best-fitting mass). The left panel of Fig. 5 shows the contour plot for the TF model. Our best fit has a mass of 1.2$`{}_{0.5}{}^{+1.8}\times 10^{12}`$$`M_{}`$ and $`\beta =0.50\pm 0.4`$ ($`\chi _{\mathrm{min}}^2=`$ 25). The mass enclosed in 120 kpc is $`M(<120\mathrm{kpc})=9.0_{3.0}^{+6.0}\times 10^{11}`$$`M_{}`$ . Our results are compatible with the work of W&E99: they find a mass of $`M=1.9_{1.7}^{+3.6}\times 10^{12}`$$`M_{}`$ , even though they favour a slightly radially anisotropic velocity ellipsoid. The right panel of Fig. 5 shows the data overlaid on our best-fitting model (solid line). Visual inspection shows that the large value obtained for the minimum $`\chi ^2`$ is driven by the discrepancy between model and data in the bins at 11.5 and 33 kpc. However, at large radii our TF model with $`M=1.2\times 10^{12}`$$`M_{}`$ provides a good representation of the data. Figure 5 also shows that the favourite W&E99 model (dashed curve), having a larger mass and a more radially anisotropic velocity ellipsoid, overpredicts the Galactocentric radial velocity dispersion. On the other hand, the TF model of SCB03, for which $`M=2.5\times 10^{12}`$$`M_{}`$ and $`\beta =1.25`$, i.e. heavier halo whose ellipsoid is much more tangentially anisotropic, declines too quickly in the inner part and tends to flatten at large radii (dotted curve), not following the trend shown by the data. The comparison of the fits produced by the constant anisotropy TF and NFW models shows that the latter reproduces better the trends in the data as a whole, from small to large radii. However, at very large radii it tends to overpredict the velocity dispersion. In this regime, the TF model provides a much better fit. This can be understood as follows. In the region between 50 and 150 kpc, where $`\sigma _{\mathrm{GSR},}`$ shows the decline, the slope of the TF model ranges between $`3`$ and $`4`$ whilst the slope of the NFW density profile is around $`2.5`$. This means that, in models with a constant velocity anisotropy, a steep dark matter density profile at large radii is favoured by the data. #### 2.3.2 Toy models for the velocity anisotropy We will now briefly relax the assumption that $`\beta `$ is constant with radius. We shall explore the following models for $`\beta (r)`$: * Model $`\beta `$-rad (Radially anisotropic). Diemand, Moore & Stadel (2004) have found in N-body $`\mathrm{\Lambda }`$CDM simulations that the anisotropy of subhalos velocities behaves as $$\beta (r)0.35\frac{r}{r_\mathrm{v}},\text{for}rr_\mathrm{v}.$$ (12) We will use this cosmologically motivated functional form to study the effect of an increasingly radially anisotropic velocity ellipsoid in our modelling of the radial velocity dispersion curve. * Model $`\beta `$-tg (Tangentially anisotropic). Proper motion measurements of the Magellanic Clouds and Sculptor, Ursa Minor, and Fornax dwarf spheroidals suggest that the tangential velocities of these objects are larger than their radial motions (Kroupa & Bastian 1997; Schweitzer et al. 1995, 1997; Dinescu et al. 2004). If confirmed, this would have as consequence that the velocity ellipsoid should be tangentially anisotropic at large radii. To explore the effect on our dynamical models of a velocity ellipsoid that becomes increasingly more tangential, we consider the following toy-model: $$\beta (r)=\beta _0\frac{r^2}{h^2},$$ (13) where we set the scale factor $`h=`$ 120 kpc. We choose two values for $`\beta _0`$: in the first case we arbitrarily fix it to 1 (model $`\beta `$-tg<sub>toy</sub>); in the second model ($`\beta `$-tg<sub>SN</sub>) we use a a sample of 91 nearby halo stars from Beers et al. (2000) – within 0.5 kpc from the Sun and with \[Fe/H\]$`<`$1.5 – to normalize our model. In this case, we find that $`\beta (R_{})=0.33`$ and therefore, $`\beta _0=0.33+R_{}^2/h^2`$. Using the models for $`\beta (r)`$ described above, we perform again the $`\chi ^2`$ best-fitting procedure for an NFW model of $`c=18`$. There is, therefore, in all cases, only one free parameter: the virial mass. The results of this new analysis are shown in the bottom panel of Fig. 6. The $`\beta `$-$`rad`$ model, for which the velocity ellipsoid becomes more radially anisotropic with radius, has $`\chi _{\mathrm{min}}^2=15`$. Even though the predicted radial velocity dispersion of this model does decrease with radius, this decline is of insufficient amplitude to reproduce the trend shown by the data. Note that this model, motivated by dark-matter simulations, provides a poorer fit than the constant $`\beta `$ model. Models where the velocity ellipsoid becomes more tangentially anisotropic with radius, $`\beta `$-tg<sub>toy</sub> and $`\beta `$-tg<sub>SN</sub>, follow very well the data, and have $`\chi _{\mathrm{min}}^2=6`$ and 7, respectively. We find that, for model $`\beta `$-tg<sub>toy</sub>, the best-fitting virial mass is $`M_\mathrm{v}=`$ 8.8 ($`\pm `$ 0.7, $`\pm `$ 1.2)$`\times 10^{11}`$$`M_{}`$ (at the 1$`\sigma `$, 2$`\sigma `$ level), and $`M_\mathrm{v}=`$1.5 ($`\pm `$ 0.1, $`\pm `$ 0.2)$`\times 10^{12}`$$`M_{}`$ (at the 1$`\sigma `$, 2$`\sigma `$ level) for $`\beta `$-tg<sub>SN</sub>. For the $`\beta `$-tg<sub>toy</sub> model, we find that mass enclosed in 120 kpc is $`M(<120\mathrm{kpc})=5.9\pm 0.5\times 10^{11}`$$`M_{}`$ ; for the $`\beta `$-tg<sub>SN</sub> model, $`M(<120\mathrm{kpc})=9.0\pm 0.6\times 10^{11}`$$`M_{}`$ . Table 2 summarizes the best-fitting parameters for our favourite models. This analysis highlights the mass-anisotropy degeneracy, since it shows that, even for the same functional form of $`\beta `$, the best-fitting value of the virial mass can differ by a factor of two. Note that the best-fitting values of the virial mass for the $`\beta `$-tg<sub>toy</sub> model and the $`\beta =cst`$ are very comparable, but this is a reflection of the fact that the two anisotropy parameters are not too similar throughout a fair range of the distances probed by the sample. However, since the value of $`\beta `$ in the Solar neighbourhood is in the range 0.5$`\pm `$0.1 (Chiba & Yoshii 1998), this would tend to suggest that, given that the ellipsoid needs to be tangentially anisotropic at large radii to give a good fit to the data, a higher value of the total mass is more likely. If we apply the same kind of analysis to the pseudo-isothermal sphere mass model, it is clear that $`\beta `$ has to decrease more strongly with radius than the above $`\beta `$-tg model used in combination with the NFW profile in order to give a reasonable fit to the data. This is in line with the results of Sommer-Larsen et al. (1997). By assuming a logarithmic potential for the dark matter halo, they found a velocity ellipsoid radially anisotropic at the Solar circle ($`\beta `$ 0.5) and tangentially anisotropic for $`r`$ 20 kpc. At $`r`$ 50 kpc the expected value of $`\beta 1`$. The Sommer-Larsen et al. (1997) model is consistent with our findings out to $``$ 50 kpc; however, if we extrapolate the predicted trend for $`\beta `$ to larger Galactocentric distances, we notice that $`\beta `$ does not decrease sufficiently rapidly to explain the decline observed in our data (see also Appendix B). From the above analysis it is evident that assumptions on $`\beta `$ for a particular mass model, can strongly influence the performance of the mass model. However, not all functional forms of $`\beta `$ for a given mass model produce a good fit to the data. More accurate proper motion measurements for a larger number of halo tracers and covering a larger range in Galactocentric distances will enable us to understand which trend in radius $`\beta `$ is following and, therefore, to establish more uniquely which mass model is preferred by the data. In addition to varying the velocity anisotropy parameter $`\beta `$ as function of radius, it is also possible to consider the effect of changing the slope $`\gamma `$ of the stellar density profile of the Galactic halo. In this case, however, the data is much more restrictive in the choice of possible models, since it is well-known that $`\gamma 33.5`$ out to $`50`$ kpc (Yanny et al. 2000). Equation (2) shows that possible variations of the stellar halo power law $`\gamma `$ with radius can “conspire” with variations of $`\beta `$ to reproduce the same radial velocity dispersion profile. We examine this issue further in Appendix B. ## 3 Discussion and conclusions We have derived the radial velocity dispersion profile of the stellar halo of the Milky Way using a sample of 240 halo objects with accurate distance and radial velocity measurements. The new data from the “Spaghetti” Survey led to a significant increase in the number of known objects for Galactocentric radii beyond 50 kpc, which allowed a more reliable determination of the dispersion profile out to very large distances. Our most distant probes are located at $`120`$ kpc, which in comparison to previous works (e.g. Sommer-Larsen et al. 1997) corresponds to an increase of 70 kpc in probing the outer halo. The Galactocentric radial velocity dispersion measured is approximately constant ($`\sigma _{\mathrm{GSR},}`$ 120 km s<sup>-1</sup> ) out to 30 kpc (consistent with Ratnatunga & Freeman 1986) and then it shows a continuous decline out to the last measured point (50 $`\pm `$ 22 km s<sup>-1</sup> at 120 kpc). This fall-off has important implications for the density profile of the dark matter halo of the Milky Way. In particular, in the hypothesis of a constant velocity anisotropy, an isothermal sphere can be immediately ruled out as model for the Galactic dark halo as this predicts a nearly constant radial velocity dispersion curve. We have also considered two other possible models for the dark halo: a truncated flat (TF) and a Navarro, Frenk & White (NFW) profile. We have compared the radial velocity dispersion observed with that predicted in these models for a tracer population (stellar halo) embedded in a potential provided by the dark halo. By means of a $`\chi ^2`$ test, we were able to derive the characteristic parameters and velocity anisotropy of these models that are most consistent with the observed data. In the case of a TF profile, the favourite model for the Milky Way dark matter halo has a mass $`M=1.2_{0.5}^{+1.8}\times 10^{12}`$$`M_{}`$ , with a corresponding velocity anisotropy $`\beta =0.50\pm 0.4`$. The data are also compatible with an NFW dark halo of $`M_\mathrm{v}=0.8_{0.2}^{+1.2}\times 10^{12}`$ $`M_{}`$ and $``$0.3$`\beta 1`$ for a concentration $`c=`$18. The comparison of the fits produced by the constant anisotropy TF and NFW models shows that the latter reproduces better the trends in the data as a whole, from small to large radii. However, at very large radii it tends to overpredict the velocity dispersion. In this regime the TF model –having a steeper density profile– provides a much better fit. Our determination of the dark halo mass of the Milky Way is consistent with previous works: the preferred TF model of W&E99 gives a mass $`M=1.9\times 10^{12}`$$`M_{}`$ , with a 1-$`\sigma `$ range of $`0.2<M[10^{12}`$$`M_{}`$$`]<5.5`$ and $`0.4<\beta <0.7`$; the favourite model from Klypin et al. 2002 gives $`M=1.0\times 10^{12}`$ $`M_{}`$ with $`c=`$12. However, the radial velocity dispersion predicted by these two models is larger than the observed one. The discrepancy between the observed low values of the radial velocity dispersion at large radii and that predicted for heavy dark halos raises the question of whether the velocity dispersion in the two most distant bins may be affected by systematics, such as the presence of streams, which could lower their values. The two bins in question are centered at $`90`$ kpc and $`120`$ kpc, and contain 6 and 3 objects respectively: 4 satellite galaxies and 5 globular clusters. The minimum angular separation of any two objects in these bins is 40, for the satellites, and 49, for the GCs. When considering the sample with 9 objects only two of these objects appear to be close on the sky: one globular cluster and one satellite galaxy that are located at $`(l,b)`$ (241,42). Although these are at similar distances of 96 kpc and 89 kpc, respectively, their line of sight radial velocities differ by more than 140 km s<sup>-1</sup> , thus making any physical association extremely unlikely. We have also investigated the effect of a velocity anisotropy that varies with radius on the velocity dispersion $`\sigma _{\mathrm{GSR},}`$ in the case of an NFW halo of concentration $`c=`$18. We find that the velocity anisotropy, which is radial at the Solar neighborhood, needs to become more tangentially anisotropic with radius in order to fit the observed rapid decline in $`\sigma _{\mathrm{GSR},}`$. In the case of an isothermal dark matter halo, the $`\beta `$ profile needs to decline even more steeply than in the NFW case in order to fit the data. We conclude that the behaviour of the observed velocity dispersion can be explained either by a dark matter halo following a steep density profile at large radii and constant velocity anisotropy, or by a halo with a less steep profile whose velocity ellipsoid becomes tangentially anisotropic at large radii. In order to distinguish between an NFW profile and a TF model, proper motions are fundamental since they enable the direct determination of the velocity anisotropy profile. Proper motions of GCs and satellites are becoming available (Dinescu et al. 1999; Piatek et al. 2003; Dinescu et al. 2004) albeit with large errors because of the very distant location of these objects. We may have to wait until Gaia is launched to determine the density profile of the Galactic dark-matter halo. ## Acknowledgments Giuseppina Battaglia gratefully acknowledges Eline Tolstoy. We thank Simon White and the anonymous referee for useful suggestions. This research has been partially supported by the Netherlands Organization for Scientific Research (NWO), and the Netherlands Research School for Astronomy (NOVA). Heather Morrison, Edward W. Olszewski and Mario Mateo acknowledge support from the NSF on grants AST 96-19490, AST 00-98435, AST 96-19524, AST 00-98518, AST 95-28367, AST 96-19632 and AST 00-98661. We used the web catalog http://physwww.mcmaster.ca/$``$harris/mwgc.dat for the Milky Way globular cluster data. ## Appendix A The Galactocentric radial velocity $`v_{\mathrm{GSR}}`$ (i.e. the l.o.s. heliocentric velocity $`V_{\mathrm{los}}`$ corrected for the solar motion and LSR motion) is related to the true Galactocentric radial, $`v_r`$, and tangential, $`v_\mathrm{t}`$, velocity by $$v_{\mathrm{GSR}}=v_r\widehat{ϵ_r}\widehat{ϵ_\mathrm{R}}+v{}_{\mathrm{t}}{}^{}\widehat{ϵ_\mathrm{t}}\widehat{ϵ_\mathrm{R}}$$ (14) where $`\widehat{ϵ_r}`$ is the unit vector in the radial direction towards the object as seen from the Galactic centre, $`\widehat{ϵ_\mathrm{t}}`$ is the unit vector in tangential direction in the same reference frame, and $`\widehat{ϵ_\mathrm{R}}`$ is the unit vector in the radial direction from the Sun to the object. The two scalar products depend on the heliocentric and galactocentric distances ($`d`$ and $`r`$) and position on the sky of the object ($`\varphi `$, $`\theta `$). For a given distribution function $`f(\overline{r},\overline{v})`$, the velocity dispersion profile (seen from the Sun) is given by $`\sqrt{v_{\mathrm{GSR}}^2}`$, and can be found by squaring Eq. (15) and integrating over all the velocities and averaging over the solid angle: $`v_{\mathrm{GSR}}^2\|_{\mathrm{\Omega }\mathrm{av}}={\displaystyle \frac{1}{{\displaystyle d^2\mathrm{\Omega }}}}[{\displaystyle }d^2\mathrm{\Omega }k(r,\theta ,\varphi ){\displaystyle }d^3vv_r^2f(\overline{r},\overline{v})+`$ $`{\displaystyle }d^2\mathrm{\Omega }h(r,\theta ,\varphi ){\displaystyle }d^3vv_\mathrm{t}^2f(\overline{r},\overline{v})],`$ or $`v_{\mathrm{GSR}}^2\|_{\mathrm{\Omega }\mathrm{av}}={\displaystyle \frac{1}{4\pi }}({\displaystyle }d^2\mathrm{\Omega }k(r,\theta ,\varphi )v_r^2+`$ $`+{\displaystyle }d^2\mathrm{\Omega }h(r,\theta ,\varphi )v_\mathrm{t}^2),`$ (15) where we have defined $$\widehat{ϵ_\mathrm{R}}=\frac{\overline{r}\overline{R}_{}}{d},$$ $$k(r,\theta ,\varphi )=(\widehat{ϵ_r}\widehat{ϵ_\mathrm{R}})^2=\left(\frac{r+R_{}\mathrm{cos}\varphi \mathrm{sin}\theta }{d}\right)^2,$$ and $$h(r,\theta ,\varphi )=(\widehat{ϵ_\mathrm{t}}\widehat{ϵ_\mathrm{R}})^2=\frac{R_{}^{}{}_{}{}^{2}}{d^2}(\mathrm{cos}^2\theta \mathrm{cos}^2\varphi +\mathrm{sin}^2\varphi ).$$ Eq. (15) can thus be expressed as $$v_{\mathrm{GSR}}^2\|_{\mathrm{\Omega }\mathrm{av}}=v_r^2K(r)+v_\mathrm{t}^2H(r).$$ If we assume that $`v_\theta ^2=v_\varphi ^2`$, and from the definition of the velocity anisotropy $`\beta `$ we find $`v_\mathrm{t}^2=2v_r^2(1\beta )`$, then it follows that $$v_{\mathrm{GSR}}^2\|_{\mathrm{\Omega }\mathrm{av}}=v_{r}^{}{}_{}{}^{2}[K(r)+2(1\beta )H(r)].$$ (16) By assuming $`v_r=0`$ and $`v{}_{\mathrm{t}}{}^{}=0`$, it follows that $`v{}_{\mathrm{GSR}}{}^{}=0`$; by performing the above integrals for $`r>R_{}`$, we find that the Galactocentric radial velocity dispersion is related to the true radial velocity dispersion by $$\sigma {}_{\mathrm{GSR}}{}^{}(r)=\sigma _r(r)\sqrt{1+2(1\beta )H(r)},$$ (17) where $$H(r)=\frac{r^2+R_{}^2}{4r^2}\frac{(r^2R_{}^2)^2}{8r^3R_{}}\mathrm{ln}\frac{r+R_{}}{rR_{}}.$$ (18) ## Appendix B Equation (2) shows that the radial velocity dispersion profile depends on the circular velocity given by the dominant mass component (i.e. the dark matter halo), the velocity anisotropy parameter $`\beta `$ and the power $`\gamma `$ of the density profile of the tracer population. For constant $`\beta `$ and $`\gamma `$, we can rewrite Eq.(2) as $$\sigma _{}^{2}{}_{r,}{}^{}(r)=\frac{1}{r^{2\beta \gamma }}_{r}^{}{}_{}{}^{\mathrm{}}V_\mathrm{c}^2(r^{})r_{}^{}{}_{}{}^{\mathrm{\hspace{0.17em}2}\beta \gamma 1}𝑑r^{}.$$ (19) In our work we assumed $`\gamma =`$ 3.5 at all Galactocentric distances, but the above equation shows also that for a fixed mass distribution (i.e. fixed circular velocity), models with the same value for $`2\beta \gamma `$ give rise to the same radial velocity dispersion profile. In this Section we explore how $`\beta `$ or $`\gamma `$ have to vary together in order to reproduce the observed Galactocentric radial velocity dispersion. In this analysis we restrict ourselves to Galactocentric distances larger than 40 kpc, where: the value of $`\gamma `$ starts to become more uncertain, the observed Galactocentric radial velocity dispersion declines and the correction factor between the Galactocentric and the true radial velocity dispersions is negligible. At these distances the Galactocentric radial velocity dispersion profile is well represented by a straight line, $`\sigma _{\mathrm{GSR},\mathrm{fit}}=ar+b`$, with $`a=`$0.6 and $`b=`$ 132 (Fig. 7, left). We assume that the circular velocity for the dark matter halo is constant and we fix it to $`V_\mathrm{c}(r)=V_\mathrm{c}=`$ 220 km s<sup>-1</sup> . By solving the Eq.(2) we obtain $$\sigma _{r,}^2=\frac{V_\mathrm{c}^2}{\gamma 2\beta }$$ (20) For all the values of $`\beta `$ and $`\gamma `$ that satify this relation (at every $`r`$) the predicted radial velocity dispersion curve will be the same. By imposing $`\sigma _{r,}^2=\sigma _{\mathrm{GSR},\mathrm{fit}}^2`$ in Eq.(19), it follows $$\gamma 2\beta =\frac{V_{\mathrm{c}}^{}{}_{}{}^{2}}{\sigma _{\mathrm{GSR},\mathrm{fit}}^2}$$ (21) Figure 7 (second panel) shows the above relation for the assumed model. The third panel in Fig. 7 shows how $`\beta `$ has to vary with the Galactocentric distance for this model if we fix $`\gamma =`$3.5, whilst the panel on the right shows how $`\gamma `$ has to change if we use the $`\beta `$-tg<sub>SN</sub> model for $`\beta `$. Clearly for this model the values the $`\gamma `$ should assume in order to reproduce the data are unrealistic. The same kind of relation between $`\beta `$ and $`\gamma `$ can be derived for different circular velocities in the regime where they can be approximated by power-laws.
warning/0506/hep-th0506236.html	ar5iv	text	# Untitled Document UTTG-01-05 Quantum Contributions to Cosmological Correlations Steven Weinberg<sup></sup><sup></sup>Electronic address: [email protected] Theory Group, Department of Physics, University of Texas Austin, TX, 78712 Abstract The “in-in” formalism is reviewed and extended, and applied to the calculation of higher-order Gaussian and non-Gaussian correlations in cosmology. Previous calculations of these correlations amounted to the evaluation of tree graphs in the in-in formalism; here we also consider loop graphs. It turns out that for some though not all theories, the contributions of loop graphs as well as tree graphs depend only on the behavior of the inflaton potential near the time of horizon exit. A sample one-loop calculation is presented. I. INTRODUCTION The departures from cosmological homogeneity and isotropy observed in the cosmic microwave background and large scale structure are small, so it is natural that they should be dominated by a Gaussian probability distribution, with bilinear averages given by the terms in the Lagrangian that are quadratic in perturbations. Nevertheless, there is growing interest in the possibility of observing non-Gaussian terms in various correlation functions,<sup>1</sup> such as an expectation value of a product of three temperature fluctuations. It is also important to understand the higher-order corrections to bilinear correlation functions, which appear in Gaussian correlations. Until now, higher-order cosmological correlations have been calculated by solving the classical field equations beyond the linear approximation. As will be shown in the Appendix, this is equivalent to calculating sums of tree graphs, though in a formalism different from the familiar Feynman graph formalism. For instance, Maldacena<sup>2</sup> has calculated the non-Gaussian average of a product of three scalar and/or gravitational fields to first order in their interactions, which amounts to calculating a tree graph consisting of a single vertex with 3 attached gravitational and/or scalar field lines. This paper will discuss how calculations of cosmological correlations can be carried to arbitrary orders of perturbation theory, including the quantum effects represented by loop graphs. So far, loop corrections to correlation functions appear to be much too small ever to be observed. The present work is motivated by the opinion that we ought to understand what our theories entail, even where in practice its predictions cannot be verified experimentally, just as field theorists in the 1940s and 1950s took pains to understand quantum electrodynamics to all orders of perturbation theory, even though it was only possible to verify results in the first few orders. There is a particular question that will concern us. In the familiar calculations of lowest-order Gaussian correlations, and also in Maldacena’s tree-graph calculation of non-Gaussian correlations, the results depended only on the behavior of the unperturbed inflaton field near the time of horizon exit. Is the same true for loop graphs? If so, it will be possible to calculated the loop contributions with some confidence, but we can learn little new from such calculations. On the other hand, if the contribution of loop graphs depends on the whole history of the unperturbed inflaton field, then calculations become much more difficult, but potentially more revealing. In this case, it might even be that the loop contributions are much larger than otherwise expected. The appropriate formalism for dealing with this sort of problem is the “in-in” formalism originally due to Schwinger.<sup>3</sup> Schwinger’s presentation is somewhat opaque, so this formalism is outlined (and extended) in an Appendix. In section II we summarize those aspects of this formalism that are needed for our present purposes. Section III introduces a class of theories to serve as a basis of discussion, with a single inflaton field, plus any number of additional massless scalar fields with only gravitational interactions and vanishing expectation values. In Section IV we prove a general theorem about the late time behavior of cosmological correlations at fixed internal as well as external wave numbers. Section V introduces a class of unrealistic theories to illustrate the problems raised by the integration over internal wave numbers, and how these problems may be circumvented. In Section VI we return to the theories introduced in Section III, and we show that the conditions of the theorem proved in Section IV are satisfied for these theories. This means that, to all orders of perturbation theory, if ultraviolet divergences cancel in the integrals over internal wave numbers, then cosmological correlations do indeed depend only on the behavior of the unperturbed inflation field near the time of horizon exit in the cases studied. We can also find other theories in which this result does not apply, as for instance by giving the additional scalar fields a self-interaction. Section VII presents a sample one-loop calculation of a cosmological correlation. II. THE “IN-IN” FORMALISM IN COSMOLOGY The problem of calculating cosmological correlation functions differs from the more familiar problems encountered in quantum field theory in at least three respects: We are not interested here in the calculation of S-matrix elements, but rather in evaluating expectation values of products of fields at a fixed time. * Conditions are not imposed on the fields at both very early and very late times, as in the calculation of S-matrix elements, but only at very early times, when the wavelength is deep inside the horizon and according to the Equivalence Principle the interaction picture fields should have the same form (when expressed in terms of metric rather than co-moving coordinates) as in Minkowski spacetime. * Although the Hamiltonian $`H`$ that generates the time dependence of the various quantum fields is constant in time, the time-dependence of the fluctuations in these fields are governed by a fluctuation Hamiltonian $`\stackrel{~}{H}`$ with an explicit time dependence, which as shown in the Appendix is constructed by expanding $`H`$ around the unperturbed solution of the field equation, and discarding the terms of first order in the perturbations to the fields and their canonical conjugates. Given a fluctuation Hamiltonian $`\stackrel{~}{H}`$, we want to use it to calculate expectation values of some product $`Q(t)`$ of field operators, all at the same time $`t`$ but generally with different space arguments. As discussed in the Appendix, the prescription of the “in-in” formalism is that $$Q(t)=\left[\overline{T}\mathrm{exp}\left(i_{\mathrm{}}^tH_I(t)𝑑t\right)\right]Q^I(t)\left[T\mathrm{exp}\left(i_{\mathrm{}}^tH_I(t)𝑑t\right)\right],$$ (1) Here $`T`$ denotes a time-ordered product; $`\overline{T}`$ is an anti-time-ordered product; $`Q^I`$ is the product $`Q`$ in the interaction picture (with time-dependence generated by the part of $`\stackrel{~}{H}`$ that is quadratic in fluctuations); and $`H_I`$ is the interaction part of $`\stackrel{~}{H}`$ in the interaction picture. (This result is different from that originally given by Maldacena<sup>2</sup> and other authors<sup>4</sup>, who left out the time-ordering and anti-time-ordering, perhaps through a typographical error. However, this makes no difference to first order in the interaction, which is the approximation used by these authors in their calculations.) We are here taking the time $`t_0`$ at which the fluctuations are supposed to behave like free fields as $`t_0=\mathrm{}`$, which is appropriate for cosmology because at very early times the fluctuation wavelengths are deep inside the horizon. Eq. (1) leads to a fairly complicated diagrammatic formalism, described in the Appendix. Unfortunately this formalism obscures crucial cancellations that occur between different diagrams. For our present purposes, it is more convenient to use a formula equivalent to Eq. (1): $`Q(t)`$ $`=`$ $`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}i^N{\displaystyle _{\mathrm{}}^t}𝑑t_N{\displaystyle _{\mathrm{}}^{t_N}}𝑑t_{N1}\mathrm{}{\displaystyle _{\mathrm{}}^{t_2}}𝑑t_1`$ (2) $`\times [H_I(t_1),[H_I(t_2),\mathrm{}[H_I(t_N),Q^I(t)]\mathrm{}]],`$ (with the $`N=0`$ term understood to be just $`Q^I(t)`$). This can easily be derived from Eq. (1) by mathematical induction. Obviously Eqs. (1) and (2) give the same results to zeroth and first order in $`H_I`$. If we assume that the right-hand sides of Eqs. (1) and (2) are equal for arbitrary operators $`Q`$ up to order $`N1`$ in $`H_I`$, then by differentiating these equations we easily see that the time derivatives of the right-hand sides are equal up to order $`N`$. Eqs. (1) and (2) also give the same results for $`t\mathrm{}`$ to all orders, so they give the same results for arbitrary $`t`$ to order $`N`$. III. THEORIES OF INFLATION To make our discussion concrete, in this section we will take up a particular class of theories of inflation. The reader who prefers to avoid details of specific theories can skip this section, and go on immediately to the general analysis of late-time behavior in the following section. In this section we will consider theories of inflation with two kinds of matter fields : a real scalar field $`\phi (𝐱,t)`$ with a non-zero homogeneous expectation value $`\overline{\phi }(t)`$ that rolls down a potential $`V(\phi )`$, and any number of real massless scalar fields $`\sigma _n(𝐱,t)`$, which have only minimal gravitational interactions, and are prevented by unbroken symmetries from acquiring vacuum expectation values. The real field $`\phi `$ serves as an inflaton whose energy density drives inflation, while the $`\sigma _n`$ are a stand-in for the large number of species of matter fields that will dominate the effects of loop graphs on the correlations of the inflaton field.<sup>\</sup><sup>\</sup>\Standard counting arguments show that in these theories the number of factors of $`8\pi G`$ in any graph equals the number of loops of any kind, plus a fixed number that depends only on which correlation function is being calculated. Matter loops are numerically more important than loops containing graviton or inflaton lines, because they carry an additional factor equal to the number of types of matter fields. We follow Maldacena,<sup>2</sup> adopting a gauge in which there are no fluctuations in the inflaton field, so that $`\phi (𝐱,t)=\overline{\phi }(t)`$, and in which the spatial part of the metric takes the form<sup>\</sup><sup>\</sup>\I am adopting Maldacena’s notation, but the quantity he calls $`\zeta `$ is more usually called $``$. To first order in fields, the quantity usually called $`\zeta `$ is defined as $`\mathrm{\Psi }H\delta \rho /\dot{\overline{\rho }}`$, while the quantity usually called $``$ is defined as $`\mathrm{\Psi }+H\delta u`$. (Here the contribution of scalar modes to $`g_{ij}`$ is written in general gauges as $`2a^2(\mathrm{\Psi }\delta _{ij}+^2\mathrm{\Psi }^{}/x^ix^j)`$, while $`\delta \rho `$ and $`\overline{\rho }`$ are the perturbation to the total energy density and its unperturbed value, while $`\delta u`$ is the perturbed velocity potential, which for a single inflaton field is $`\delta u=\delta \phi /\dot{\overline{\phi }}`$.) In the gauge used by Maldacena and in the present paper $`\delta u=\mathrm{\Psi }^{}=0`$, so since $`\zeta `$ is defined here as $`\mathrm{\Psi }`$ to first order in fields, it corresponds to the quantity usually called $``$. Outside the horizon $``$ and $`\zeta `$ are the same. $$g_{ij}=a^2e^{2\zeta }[\mathrm{exp}\gamma ]_{ij},\gamma _{ii}=0,_i\gamma _{ij}=0.$$ (3) where $`a(t)`$ is the Robertson–Walker scale factor, $`\gamma _{ij}(𝐱,t)`$ is a gravitational wave amplitude, and $`\zeta (𝐱,t)`$ is a scalar whose characteristic feature is that it is conserved outside the horizon,<sup>5</sup> that is, for physical wave numbers that are small compared with the expansion rate. The same is true of $`\gamma _{ij}`$. The other components of the metric are given in the Arnowitt–Deser–Misner (ADM) formalism<sup>6</sup> by $$g_{00}=N^2+g_{ij}N^iN^j,g_{i0}=g_{ij}N^j,$$ (4) where $`N`$ and $`N^i`$ are auxiliary fields, whose time-derivatives do not appear in the action. The Lagrangian density in this gauge (with $`8\pi G1`$) is $`={\displaystyle \frac{a^3}{2}}e^{3\zeta }[NR^{(3)}2NV(\overline{\phi })+N^1\left(E^j{}_{i}{}^{}E_{}^{i}{}_{j}{}^{}(E^i{}_{i}{}^{})^2\right)+N^1\dot{\overline{\phi }}^2`$ $`+N^1{\displaystyle \underset{n}{}}(\dot{\sigma }_nN^i_i\sigma _n)^2Na^2e^{2\zeta }[\mathrm{exp}(\gamma )]^{ij}{\displaystyle \underset{n}{}}_i\sigma _n_j\sigma _n],`$ where $$E_{ij}\frac{1}{2}\left(\dot{g}_{ij}_iN_j_jN_i\right),$$ (6) and bars denote unperturbed quantities. All spatial indices $`i`$, $`j`$, etc. are lowered and raised with the matrix $`g_{ij}`$ and its reciprocal; $`_i`$ is the three-dimensional covariant derivative calculated with this three-metric; and $`R^{(3)}`$ is the curvature scalar calculated with this three-metric: $$R^{(3)}=a^2e^{2\zeta }\left[e^\gamma \right]^{ij}R_{ij}^{(3)}.$$ The auxiliary fields $`N`$ and $`N^i`$ are to be found by requiring that the action is stationary in these variables. This gives the constraint equations: $$_i\left[N^1(E^i{}_{j}{}^{}\delta ^i{}_{j}{}^{}E_{}^{k}{}_{k}{}^{})\right]=N^1\underset{n}{}_j\sigma _n(\dot{\sigma }_nN^i_i\sigma _n),$$ (7) $`N^2[R^{(3)}2Va^2e^{2\zeta }[\mathrm{exp}(\gamma )]^{ij}{\displaystyle \underset{n}{}}_i\sigma _n_j\sigma _n]=E^i{}_{j}{}^{}E_{}^{j}{}_{i}{}^{}(E^i{}_{i}{}^{})^2`$ $`+\dot{\overline{\phi }}^2+{\displaystyle \underset{n}{}}\left(\dot{\sigma }_nN^i_i\sigma _n\right)^2`$ (8) For instance, to first order in fields (including field derivatives) the auxiliary fields are the same as in the case of no additional matter fields<sup>2</sup> $$N=1+\dot{\zeta }/H,N^i=\frac{1}{a^2H}_i\zeta +ϵ_i^2\dot{\zeta },$$ (9) where $$ϵ\frac{\dot{H}}{H^2}=\frac{\dot{\overline{\phi }}^2}{2H^2},H\frac{\dot{a}}{a}$$ (10) The fields in the interaction picture satisfy free-field equations. For $`\zeta `$ we have the Mukhanov equation:<sup>7</sup> $$\frac{^2\zeta }{t^2}+\left[\frac{d\mathrm{ln}(a^3ϵ)}{dt}\right]\frac{\zeta }{t}a^2^2\zeta =0,$$ (11) The field equation for gravitational waves is $$\frac{^2\gamma _{ij}}{t^2}+3H\frac{\gamma _{ij}}{t}a^2^2\gamma _{ij}=0,$$ (12) and for the matter fields $$\frac{^2\sigma _n}{t^2}+3H\frac{\sigma _n}{t}a^2^2\sigma _n=0.$$ (13) The fields in the interaction picture are then $$\zeta (𝐱,t)=d^3q\left[e^{i𝐪𝐱}\alpha (𝐪)\zeta _q(t)+e^{i𝐪𝐱}\alpha ^{}(𝐪)\zeta _q^{}(t)\right],$$ (14) $$\gamma _{ij}(𝐱,t)=d^3q\underset{\lambda }{}\left[e^{i𝐪𝐱}e_{ij}(\widehat{q},\lambda )\alpha (𝐪,\lambda )\gamma _q(t)+e^{i𝐪𝐱}e_{ij}^{}(\widehat{q},\lambda )\alpha ^{}(𝐪,\lambda )\gamma _q^{}(t)\right],$$ (15) $$\sigma _n(𝐱,t)=d^3q\left[e^{i𝐪𝐱}\alpha (𝐪,n)\sigma _q(t)+e^{i𝐪𝐱}\alpha ^{}(𝐪,n)\sigma _q^{}(t)\right],$$ (16) where $`\lambda =\pm 2`$ is a helicity index and $`e_{ij}(\widehat{q},\lambda )`$ is a polarization tensor, while $`\alpha (𝐪)`$, $`\alpha (𝐪,\lambda )`$, and $`\alpha (𝐪,n)`$ are conventionally normalized annihilation operators, satisfying the usual commutation relations $$[\alpha (𝐪),\alpha ^{}(𝐪^{})]=\delta ^3\left(𝐪𝐪^{}\right),[\alpha (𝐪),\alpha (𝐪^{})]=0.$$ (17) $$[\alpha (𝐪,\lambda ),\alpha ^{}(𝐪^{},\lambda ^{})]=\delta _{\lambda \lambda ^{}}\delta ^3\left(𝐪𝐪^{}\right),[\alpha (𝐪,\lambda ),\alpha (𝐪^{},\lambda ^{})]=0,$$ (18) and $$[\alpha (𝐪,n),\alpha ^{}(𝐪^{},n^{})]=\delta _{nn^{}}\delta ^3\left(𝐪𝐪^{}\right),[\alpha (𝐪,n),\alpha (𝐪^{},n^{})]=0,$$ (19) Also, $`\zeta _q(t)`$, $`\gamma _q(t)`$, and $`\sigma _q(t)`$ are suitably normalized positive-frequency solutions of Eqs. (11)–(13), with $`^2`$ replaced with $`q^2`$. They satisfy initial conditions, designed to make $`\zeta \dot{\overline{\phi }}/H`$, $`\gamma _{ij}/\sqrt{16\pi G}`$, and $`\sigma _n`$ behave like conventionally normalized free fields at $`t\mathrm{}`$:In Newtonian gauge the quantity $`\zeta (𝐱,t)\dot{\overline{\phi }}(t)/H(t)`$ approaches the inflaton field fluctuation $`\delta \phi (t)`$ for $`t\mathrm{}`$. $`{\displaystyle \frac{\dot{\overline{\phi }}(t)\zeta _q(t)}{H(t)}}{\displaystyle \frac{\gamma _q(t)}{\sqrt{16\pi G}}}\sigma _q(t)`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}\sqrt{2q}a(t)}}\mathrm{exp}\left(iq{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{dt^{}}{a(t^{})}}\right).`$ (20) IV. LATE TIME BEHAVIOR The question to be addressed in this section is whether the time integrals in Eqs. (1) and (2) are dominated by times near horizon exit for general graphs. This question is more complicated for loop graphs than for tree graphs, such as that considered by Maldacena, because for loops there are two different kinds of wave number: the fixed wave numbers $`q`$ associated with external lines, and the internal wave numbers $`p`$ circulating in loops, over which we must integrate. It is only if the integrals over internal wave numbers $`p`$ are dominated by values of order $`pq`$ that we can speak of a definite time of horizon exit, when $`q/ap/aH`$. In this section we will integrate first over the time arguments in Eq. (2), holding the internal wave numbers at fixed values, and return at the end of this section to the problems raised by the necessity of then integrating over the $`p`$s. There is never any problem with the convergence of the time integrals at very early times; all fluctuations oscillate very rapidly for $`q/aH`$ and $`p/aH`$, suppressing the contribution of early times to the time integrals in Eq. (2). To see what happens for late times, when $`q/aH`$ and $`p/aH`$, we need to count the powers of $`a`$ in the contribution of late times in general loop as well as tree graphs. For this purpose, we need to consider the behavior of the coefficient functions appearing in the Fourier decompositions (14)–(16) of the fields in the interaction picture. In order to implement dimensional regularization, we will consider these coefficient functions in $`2\nu `$ space dimensions, returning later to the limit $`2\nu 3`$. The coefficient functions then obey differential equations obtained by replacing the space dimensionality 3 in Eqs. (11)–(13) with $`2\nu `$, as well as replacing the Laplacian with $`q^2`$: $$\frac{d^2\zeta _q(t)}{dt^2}+\left[\frac{d\mathrm{ln}\left(a^{2\nu }(t)ϵ(t)\right)}{dt}\right]\frac{d\zeta _q(t)}{dt}+\frac{q^2}{a^2(t)}\zeta _q(t)=0,$$ (21) $$\frac{d^2\gamma _q(t)}{dt^2}+2\nu H(t)\frac{d\gamma _q(t)}{dt}+\frac{q^2}{a^2(t)}\gamma _q(t)=0,$$ (22) $$\frac{d^2\sigma _q(t)}{dt^2}+2\nu H(t)\frac{d\sigma _q(t)}{dt}+\frac{q^2}{a^2}\sigma _q(t)=0.$$ (23) At late times, when $`q/aH`$, the solutions can be written as asymptotic expansions in inverse powers of $`a`$:<sup>††</sup><sup>††</sup>††By $`t=\mathrm{}`$ in the limits of these integrals and elsewhere in this paper, we mean a time still during inflation, but sufficiently late so that $`a(t)`$ is many e-foldings larger than its value when $`q/a`$ falls below $`H`$. $`\zeta _q(t)\zeta _q^o\left[1+{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{q^2dt^{}}{a^{2\nu }(t^{})ϵ(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a^{2\nu 2}(t^{\prime \prime })ϵ(t^{\prime \prime })𝑑t^{\prime \prime }+\mathrm{}\right]`$ $`+C_q[{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})ϵ(t^{})}}`$ $`+q^2{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})ϵ(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a^{2\nu 2}(t^{\prime \prime })ϵ(t^{\prime \prime })dt^{\prime \prime }{\displaystyle _{t^{\prime \prime }}^{\mathrm{}}}{\displaystyle \frac{dt^{\prime \prime \prime }}{a^{2\nu }(t^{\prime \prime \prime })ϵ(t^{\prime \prime \prime })}}+\mathrm{}]`$ (24) $`\gamma _q(t)\gamma _q^o\left[1+{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{q^2dt^{}}{a^{2\nu }(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a^{2\nu 2}(t^{\prime \prime })𝑑t^{\prime \prime }+\mathrm{}\right]`$ $`+D_q\left[{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})}}+q^2{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a(t^{\prime \prime })𝑑t^{\prime \prime }{\displaystyle _{t^{\prime \prime }}^{\mathrm{}}}{\displaystyle \frac{dt^{\prime \prime \prime }}{a^{2\nu }(t^{\prime \prime \prime })}}+\mathrm{}\right]`$ (25) $`\sigma _q(t)\sigma _q^o\left[1+{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{q^2dt^{}}{a^{2\nu }(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a^{2\nu 2}(t^{\prime \prime })𝑑t^{\prime \prime }+\mathrm{}\right]`$ $`+E_q\left[{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})}}+q^2{\displaystyle _t^{\mathrm{}}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a^{2\nu 2}(t^{\prime \prime })𝑑t^{\prime \prime }{\displaystyle _{t^{\prime \prime }}^{\mathrm{}}}{\displaystyle \frac{dt^{\prime \prime }}{a^{2\nu }(t^{\prime \prime })}}+\mathrm{}\right]`$ where $`\zeta _q^o`$, $`\gamma _q^0`$, and $`\sigma _q^o`$ are the limiting values of $`\zeta _q(t)`$, $`\gamma _q(t)`$, and $`\sigma _q(t)`$ (the “$`o`$” superscript stands for “outside the horizon”) and $`C_q`$, $`D_q`$, and $`E_q`$ are additional constants. In any kind of inflation with sufficient expansion, the Robertson-Walker scale factor $`a`$ will grow much faster than $`H`$ or $`ϵ`$ can change, and Eqs. (24)–(26) thus show that (at least for $`2\nu 2`$) the time derivatives of $`\zeta _q`$, $`\gamma _q`$, and $`\sigma _q`$ all vanish for $`q/aH`$ like $`1/a^2`$. If an interaction involves enough factors of $`\dot{\zeta }`$, $`\dot{\gamma }_{ij}`$, and/or $`\dot{\sigma }_n`$ so that these $`1/a^2`$ factors and any $`1/a^2`$ factors from the contraction of space indices more than compensate for the $`a^{2\nu }`$ factor in the interaction from the square root of the metric determinant, then the integral over the associated time coordinate will converge exponentially fast at late times as well as at early times, and therefore may be expected to be dominated by the era in which the wavelength leaves the horizon. For instance, the extension of Eq. (5) to $`2\nu `$ space dimensions gives the interaction between a $`\zeta `$ field and a pair of $`\sigma `$ fields $`_{\zeta \sigma \sigma }`$ $`=`$ $`{\displaystyle \frac{a^{2\nu 2}}{2}}\zeta {\displaystyle \underset{n}{}}_i\sigma _n_i\sigma _n{\displaystyle \frac{a^{2\nu 2}}{2H}}\dot{\zeta }{\displaystyle \underset{n}{}}_i\sigma _n_i\sigma _n`$ (27) $`+a^{2\nu 2}_i\left({\displaystyle \frac{\zeta }{H}}ϵa^2^2\dot{\zeta }\right){\displaystyle \underset{n}{}}\dot{\sigma }_n_i\sigma _n`$ $`{\displaystyle \frac{a^{2\nu }}{2H}}\dot{\zeta }{\displaystyle \underset{n}{}}\dot{\sigma }_n^2+{\displaystyle \frac{3a^{2\nu }}{2}}\zeta {\displaystyle \underset{n}{}}\dot{\sigma }_n^2.`$ (The $`\zeta \sigma \sigma `$ interaction Hamiltonian given by canonical quantization is just $`d^{2\nu }x_{\zeta \sigma \sigma }`$, but this simple relation does not always apply.) Counting a factor $`a^2`$ for each $`\dot{\zeta }`$ or $`\dot{\sigma }_n`$, the terms in this interaction go as $`a^{2\nu 2}`$, $`a^{2\nu 4}`$, $`a^{2\nu 4}`$, $`a^{2\nu 6}`$, and $`a^{2\nu 4}`$, respectively. All these terms are safe for $`2\nu <4`$, except for the first, which for $`2\nu >2`$ grows exponentially at late times. Because of the commutators in Eq. (2), the condition for a safe interaction is actually less stringent than that it should decay exponentially with time, and even a growing term that only involves fields rather than their time derivatives, like the first term in Eq. (27), may not destroy the convergence of the time integrals. We will now prove the following: Theorem The integrals over the time coordinates of interactions converge exponentially for $`t\mathrm{}`$, essentially as $`^{\mathrm{}}𝑑t/a^n(t)`$ with $`n>0`$, provided that in $`2\nu `$ space dimensions, all interactions are of one or the other of two types: Safe interactions, that contain a number of factors of $`a(t)`$ (including $`2`$ factors of $`a`$ for each time derivative and the $`2\nu `$ factors of $`a`$ from $`\sqrt{\mathrm{Det}g}`$) strictly less than $`2\nu 2`$, and * Dangerous interactions, which grow at late times no faster than $`a^{2\nu 2}`$, and contain only fields, not time derivatives of fields. These conditions are evidently met by the interaction (27), irrespective of the value of $`\nu `$, and, as we shall see in Section VI, they are satisfied by all other interactions in the theories of Section III, but not in all theories. Before proceeding to the proof, it should be noted that just as in Eq. (27), the space dimensionality $`2\nu `$ enters in the interaction only in a factor $`\sqrt{\mathrm{Det}g}a^{2\nu }`$, so the question of whether or not a given theory satisfies the conditions of this theorem does not depend on the value of $`2\nu `$. Thus this theorem has the corollary: Corollary The integrals over the time coordinates of interactions converge exponentially for $`t\mathrm{}`$, essentially as $`^{\mathrm{}}𝑑t/a^n(t)`$ with $`n>0`$, provided that in 3 space dimensions all interactions are of one or the other of two types: * Safe interactions, that contain a number of factors of $`a(t)`$ (including $`2`$ factors of $`a`$ for each time derivative and the 3 factors of $`a`$ from $`\sqrt{\mathrm{Det}g}`$) strictly less than $`+1`$, and * Dangerous interactions, which grow at late times no faster than $`a`$, and contain only fields, not time derivatives of fields. Here is the proof. As already mentioned, the reason that dangerous interactions are not necessarily fatal has to do with how they enter into commutators in Eq. (2). Because of the time-ordering in Eq. (2), any failure of convergence of the time integrals for $`t+\mathrm{}`$ in $`N`$th-order perturbation theory must come from a region of the multi-time region of integration in which, for some $`r`$, the time arguments $`t_r,t_{r+1},\mathrm{}t_N`$, all go to infinity together. We will therefore have to count the number of factors of $`a(t_r),a(t_{r+1}),\mathrm{}a(t_N)`$, treating them all as being of the same order of magnitude. (This does not take proper account of factors of $`\mathrm{log}a`$, but as long as the integral over $`t_r,t_{r+1},\mathrm{}t_N`$ involves a negative total number of factors of $`a`$, it converges exponentially fast no matter how many factors of $`\mathrm{log}a`$ arise from subintegrations.) Now, at least one of the fields or field time derivatives in each term in $`H(t_s)`$ with $`rsN`$ must appear in a commutator with one of the fields in some other $`H_I(t_s^{})`$ with $`s<s^{}N`$. So we need to consider the commutators of fields at times which may be unequal, but are both late. In the sense described above, treating all $`a(t_r),a(t_{r+1}),\mathrm{}a(t_N)`$ as being of the same order of magnitude, if $`a(t)`$ increases more-or-less exponentially, then the commutator of two fields or a field and a field time-derivative goes as $`a^{2\nu }`$, while the commutator of two field time-derivatives goes as $`a^{2\nu 2}`$. For instance, the unequal-time commutators of the interaction-picture fields (14)–(16) are $$[\zeta (𝐱,t),\zeta (𝐱^{},t^{})]=d^{2\nu }pe^{i𝐩(𝐱𝐱^{})}\left(\zeta _p(t)\zeta _p^{}(t^{})\zeta _p(t^{})\zeta _p^{}(t)\right),$$ (28) $$[\gamma _{ij}(𝐱,t),\gamma _{kl}(𝐱^{},t^{})]=d^{2\nu }pe^{i𝐩(𝐱𝐱^{})}\mathrm{\Pi }_{ijkl}(\widehat{p})\left(\gamma _p(t)\gamma _p^{}(t^{})\gamma _p(t^{})\gamma _p^{}(t)\right),$$ (29) $$[\sigma _n(𝐱,t),\sigma _m(𝐱^{},t^{})]=\delta _{nm}d^{2\nu }pe^{i𝐩(𝐱𝐱^{})}\left(\sigma _p(t)\sigma _p^{}(t^{})\sigma _p(t^{})\sigma _p^{}(t)\right),$$ (30) where $`\mathrm{\Pi }_{ijkl}(\widehat{p})_\lambda e_{ij}(\widehat{p},\lambda )e_{kl}(\widehat{p},\lambda )`$. The two asymptotic expansions given in Eqs.(21–(23) for each of the fields are both real aside from over-all factors, so neither by itself contributes to the commutators. On the other hand, the constants $`C_p\zeta _p^o`$, $`D_p\gamma _p^o`$, and $`E_p\sigma _p^o`$ are in general complex. (For instance, in a strictly exponential expansion, inflation, the phase of $`C_p\zeta _p^o`$ is given by a factor $`e^{i\nu \pi }`$.) The asymptotic expansions of the commutators at late times are therefore $`[\zeta (𝐱_1,t_1),\zeta (𝐱_2,t_2)]`$ $``$ $`2i{\displaystyle }d^{2\nu }p\mathrm{Im}\left[C_p\zeta _p^o\right]e^{i𝐩(𝐱_1𝐱_2)}[{\displaystyle _{t_1}^{t_2}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})ϵ(t^{})}}`$ $`+p^2{\displaystyle _{t_1}^{t_2}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})ϵ(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a^{2\nu 2}(t^{\prime \prime })ϵ(t^{\prime \prime })𝑑t^{\prime \prime }{\displaystyle _{t^{\prime \prime }}^{\mathrm{}}}{\displaystyle \frac{dt^{\prime \prime }}{a^{2\nu }(t^{\prime \prime })ϵ(t^{\prime \prime })}}`$ $`+p^2{\displaystyle _{t_1}^{\mathrm{}}}{\displaystyle \frac{dt_1^{}}{a^{2\nu }(t_1^{})ϵ(t_1^{})}}{\displaystyle _{t_2}^{\mathrm{}}}{\displaystyle \frac{dt_2^{}}{a^{2\nu }(t_2^{})ϵ(t_2^{})}}{\displaystyle _{t_1^{}}^{t_2^{}}}a^{2\nu 2}(t^{\prime \prime })ϵ(t^{\prime \prime })dt^{\prime \prime }+\mathrm{}],`$ $`[\gamma _{ij}(𝐱_1,t_1),\gamma _{kl}(𝐱_2,t_2)]`$ $``$ $`2i{\displaystyle }d^{2\nu }p\mathrm{\Pi }_{ijkl}(\widehat{p})\mathrm{Im}\left[D_p\gamma _p^o\right]e^{i𝐩(𝐱_1𝐱_2)}[{\displaystyle _{t_1}^{t_2}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})}}`$ $`+p^2{\displaystyle _{t_1}^{t_2}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a^{2\nu 2}(t^{\prime \prime })𝑑t^{\prime \prime }{\displaystyle _{t^{\prime \prime }}^{\mathrm{}}}{\displaystyle \frac{dt^{\prime \prime }}{a^{2\nu }(t^{\prime \prime })}}`$ $`+p^2{\displaystyle _{t_1}^{\mathrm{}}}{\displaystyle \frac{dt_1^{}}{a^{2\nu }(t_1^{})}}{\displaystyle _{t_2}^{\mathrm{}}}{\displaystyle \frac{dt_2^{}}{a^{2\nu }(t_2^{})}}{\displaystyle _{t_1^{}}^{t_2^{}}}a^{2\nu 2}(t^{\prime \prime })dt^{\prime \prime }+\mathrm{}],`$ $`[\sigma _n(𝐱_1,t_1),\sigma _m(𝐱_2,t_2)]`$ $``$ $`2i\delta _{nm}{\displaystyle }d^{2\nu }p\mathrm{Im}\left[E_p\sigma _p^o\right]e^{i𝐩(𝐱_1𝐱_2)}[{\displaystyle _{t_1}^{t_2}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})}}`$ $`+p^2{\displaystyle _{t_1}^{t_2}}{\displaystyle \frac{dt^{}}{a^{2\nu }(t^{})}}{\displaystyle _{\mathrm{}}^t^{}}a^{2\nu 2}(t^{\prime \prime })𝑑t^{\prime \prime }{\displaystyle _{t^{\prime \prime }}^{\mathrm{}}}{\displaystyle \frac{dt^{\prime \prime }}{a^{2\nu }(t^{\prime \prime })}}`$ $`+p^2{\displaystyle _{t_1}^{\mathrm{}}}{\displaystyle \frac{dt_1^{}}{a^{2\nu }(t_1^{})}}{\displaystyle _{t_2}^{\mathrm{}}}{\displaystyle \frac{dt_2^{}}{a^{2\nu }(t_2^{})}}{\displaystyle _{t_1^{}}^{t_2^{}}}a^{2\nu 2}(t^{\prime \prime })dt^{\prime \prime }+\mathrm{}].`$ We see that the commutator of two fields vanishes essentially as $`a^{2\nu }`$ for late times, and the same is true for the commutator of a field and its time derivative, but the commutators of two time derivatives arise only from the third terms in the expansions (31)–(33), and therefore go as $`a^{2\nu 2}`$. That is, $`[\dot{\zeta }(𝐱_1,t_1),\dot{\zeta }(𝐱_2,t_2)]`$ $``$ $`2i{\displaystyle d^{2\nu }p\mathrm{Im}\left[C_p\zeta _p^o\right]e^{i𝐩(𝐱_1𝐱_2)}}`$ $`\times \left[{\displaystyle \frac{p^2}{a^{2\nu }(t_1)ϵ(t_1)a^{2\nu }(t_2)ϵ(t_2)}}{\displaystyle _{t_1}^{t_2}}a^{2\nu 2}(t^{})ϵ(t^{})𝑑t^{}+\mathrm{}\right],`$ and likewise for $`\gamma _{ij}`$ and $`\sigma _n`$. Let’s now add up the total number of factors of $`a(t_r),a(t_{r+1}),\mathrm{}\mathrm{and}a(t_N)`$ in the integrand of Eq. (2), for some selection of terms in the interactions $`H(t_s)`$ with $`rsN`$. Suppose that the selected term in $`H(t_s)`$ contains an explicit factor $`a(t_s)^{A_s}`$, and $`B_s`$ factors of field time derivatives. Suppose also that in the inner $`Nr+1`$ commutators in Eq. (2) there appear $`C`$ commutators of fields with each other, $`C^{}`$ commutators of fields with field time derivatives, and $`C^{\prime \prime }`$ commutators of field time derivatives with each other. The number of field time derivatives that are not in these commutators is $`_sB_sC^{}2C^{\prime \prime }`$, and these contribute a total $`2_sB_s+2C^{}+4C^{\prime \prime }`$ factors of $`a`$. (All sums over $`s`$ here run from $`r`$ to $`N`$.) In addition, there are $`_sA_s`$ factors of $`a`$ that appear explicitly in the interactions, and as we have seen, the commutators contribute $`2\nu C2\nu C^{}(2\nu +2)C^{\prime \prime }`$ factors of $`a`$. Hence the total number of factors of $`a(t_r),a(t_{r+1}),\mathrm{}\mathrm{and}a(t_N)`$ in the integrand of Eq. (2) is $$\mathrm{\#}=\underset{s}{}(A_s2B_s)2\nu C(2\nu 2)(C^{}+C^{\prime \prime })=\underset{s}{}(A_s2B_s2\nu +2)2C,$$ (34) in which we have used the fact that the total number $`C+C^{}+C^{\prime \prime }`$ of commutators of the interactions $`H(t_r),H(t_{r+1}),\mathrm{}\mathrm{and}H(t_N)`$ with each other and with the field product $`Q`$ equals the number of these interactions. Under the assumptions of this theorem, all interactions have $`A_s2B_s2\nu 2`$. If any of them are safe in the sense that $`A_s2B_s<2\nu 2`$, then $`\mathrm{\#}<0`$, and the integral over time converges exponentially fast. On the other hand, if all of them have $`A_s2B_s=2\nu 2`$, then under the assumptions of this theorem they all involve only fields, not field time derivatives, so the same is true of the commutators of these interactions. In this case $`C>0`$ and $`\mathrm{\#}=2C<0`$, so again the integral over time converges exponentially fast. In counting powers of $`a`$, we have held the wave numbers $`p`$ associated with internal lines fixed, like the external wave numbers, because we are integrating over time coordinates before we integrate over the internal wave numbers. The integrals over time receive little contribution from values of the conformal time $`\tau _t^{\mathrm{}}𝑑t/a`$ satisfying $`p\tau 1`$ and $`q\tau 1`$, because of the rapid oscillation of the integrand, and for theories satisfying the conditions of our theorem they also receive little contribution from values of $`\tau `$ with $`p\tau 1`$ and $`q\tau 1`$, because of the damping provided by negative powers of $`a`$. (Note that when $`a(t)`$ increases more or less exponentially, $`\tau `$ is of the order of $`1/aH`$.) Thus for these theories, we expect the integrals to be dominated by times for which $`1/\tau `$ is in the range from the $`q`$s to the $`ps`$s. The question then is whether the integrals over the internal wave numbers $`p`$ are dominated by values of the order of the external wave numbers $`q`$? If they are, then the results depend only on the history of inflation around the time of horizon exit, $`q\tau 1`$, or in other words, $`q/aH`$. Any integral over the internal wave numbers will in general take the form of a polynomial in the external wave numbers, with coefficients that may be divergent, plus a finite term given by a convergent integral dominated by internal wave numbers of the same order of magnitude as the fixed external wave numbers. An example of this decomposition is given in Section VII. In particular, the integral over the wave number associated with an internal line that begins and ends at the same vertex does not involve the external wave numbers, so its contribution is purely a polynomial in the wave numbers of the other lines attached to the same vertex. Just as in dealing with ultraviolet divergences in flat space quantum field theory, renormalization removes some of these ultraviolet divergent polynomial terms, and others are removed by appropriate redefinitions of the field operators. (Some examples are given in the next section.) Where redefinition of the field operators is necessary, it is only products of the redefined “renormalized” field operators whose expectation values may be expected to give results that converge at late times. If, after all such renormalizations and redefinitions, there remained ultraviolet divergences in the integrals over internal wave numbers, we could conclude that the approximation of extending the time integrals to $`+\mathrm{}`$ is not valid, and that these integrals can be taken only to some time $`t`$ late in inflation. The decrease of the integrand at wave numbers $`p`$ much greater than $`1/\tau (t)`$ would then provide the ultraviolet cut off that is still needed, but the correlation functions would exhibit the sort of time dependence that has been found in other contexts by Woodard and his collaborators,<sup>3</sup> and we would not be able to draw conclusions about correlations actually measured at times much closer to the present. The possible presence of such ultraviolet divergences that are not removed by renormalization and field redefinition is an important issue, which merits further study.<sup>†††</sup><sup>†††</sup>†††Many theories are afflicted with infrared divergences, even when $`t`$ is held fixed. The infrared divergences are attributed to the imposition of the unrealistic initial condition, that at early times all of infinite space is occupied by a Bunch–Davies vacuum. The infrared divergence can be eliminated either by taking space to be finite<sup>8</sup> or by changing the vacuum.<sup>9</sup> In any case, it is the appearance of uncancelled ultraviolet rather than infrared divergences when we integrate over internal wave numbers after taking the limit $`t\mathrm{}`$ that shows the impropriety of this interchange of limit and integral, because factors of $`1/a`$ in the integrand are typically accompanied with factors of internal wave numbers, so that the $`1/a`$ factors do not suppress the integrand for large values of $`a`$ if the integral receives contributions from arbitrarily large values of the internal wave number. But even if such ultraviolet divergences are present, it would still be possible to calculate the non-polynomial part of the integrals over internal momenta which is not ultraviolet divergent (at least in one loop order) even when the time $`t`$ is taken to infinity. Such a calculation will be presented in Section VII. V. AN EXAMPLE: EXPONENTIAL EXPANSION To clarify the issues discussed at the end of the previous section, we will examine a simple unphysical model, along with a revealing class of generalizations. First, consider a single real scalar field $`\phi (𝐱,t)`$ in a fixed de Sitter metric.This model, and much of the analysis, was suggested to me by R. Woodard, private communication. In order to implement dimensional regularization, we work in $`2\nu `$ space dimensions, letting $`\nu 3/2`$ at the end of our calculation. The Lagrangian density is taken as $$=\frac{1}{2}\sqrt{\mathrm{Detg}}g^{\mu \nu }(1+\lambda \phi ^2)_\mu \phi _\nu \phi =(1+\lambda \phi ^2)\left[\frac{a^{2\nu }}{2}\dot{\phi }^2\frac{a^{2\nu 2}}{2}(\phi )^2\right],$$ (35) where $`ae^{Ht}`$ with $`H`$ constant. (This of course can be rewritten as a free field theory, but it is instructive nonetheless, and will be generalized later in this section to interacting theories.) We follow the usual procedure of defining a canonical conjugate field $`\pi =/\dot{\phi }`$, constructing the Hamiltonian density $`=\pi \dot{\phi }`$ with $`\dot{\phi }`$ expressed in terms of $`\pi `$, dividing $``$ into a quadratic part $`_0`$ and interaction part $`_I`$, and then replacing $`\pi `$ in $`_I`$ with the interaction-picture $`\pi _I`$ given by $`\dot{\phi }=[_0/\pi ]_{\pi =\pi _I}`$. This gives an interaction $$H_I=\frac{\lambda }{2}d^{2\nu }x\left[\frac{a^2}{2}\{\frac{\phi ^2}{1+\lambda \phi ^2},\dot{\phi }^2\}+a^{2\nu 2}(\phi )^2\phi ^2\right].$$ (36) (An anticommutator is needed in the first term to satisfy the requirement that $`H_I`$ be Hermitian.) This interaction satifies the conditions of the theorem proved in the previous section for any value of the space dimensionality $`2\nu `$: the first term in the square brackets contains $`2\nu 4`$ factors of $`a`$ (counting a factor $`a^2`$ for each time derivative, so it is safe, while the second term contains $`2\nu 2`$ factors of $`a`$, and is therefore dangerous, but it only involves fields (including space derivatives), not their time derivatives, so though dangerous it still satisfies the conditions of our theorem. To first order in $`\lambda `$, the expectation value $`\phi (𝐱,t)\phi (𝐱^{},t)`$ is given by a one-loop diagram, in which a scalar field line is emitted and absorbed at the same vertex, with the two external lines also attached to this vertex. This expectation value receives contributions of three kinds: Terms in which no time derivatives act on the internal lines. This contribution is the same as would be obtained by adding effective interactions proportional to $`a^{2\nu 2}(\phi )^2`$, $`a^{2\nu 2}\phi ^2`$, or $`a^{2\nu }\dot{\phi }^2`$, all of which satisfy the conditions of the theorem of the previous section. Thus it cannot affect the conclusion that the integral over the time argument of $`H_I(t_1)`$ converges exponentially at $`t_1=+\mathrm{}`$, so that $`\phi (𝐱,t)\phi (𝐱^{},t)`$ approaches a finite limit for $`t\mathrm{}`$. Terms in which time derivatives act on both ends of the internal line. This produces an effective interaction proportional to $`a^{2\nu }\phi ^2`$, which violates the conditions of our theorem, but it can be removed by adding an $`R\phi ^2\sqrt{\mathrm{Det}g}`$ counterterm in the Lagrangian. (This cancellation is not automatic, because the condition of minimal coupling is not enforced by any symmetry.) Terms in which a time derivative acts on just one end of the internal line. This produces an effective interaction proportional to $`a^{2\nu }\phi \dot{\phi }`$, which violates the conditions of our theorem, and cannot be removed by adding a generally covariant counterterm to the Lagrangian. To see in detail what trouble is caused by the third type of contribution, note that the interaction picture scalar field is given by a Fourier decomposition like Eq. (16), with coefficient functions<sup>‡‡</sup><sup>‡‡</sup>‡‡Here and below we will not be careful to extend factors like $`4\pi `$ to $`2\nu `$ space dimensions. This only affects constant terms that accompany any $`(2\nu 3)^1`$ poles. $$\phi _q(t)=\frac{e^{i\pi (2\nu +1)/4}H^{\nu 1/2}}{4\pi \sqrt{2}q^\nu }H_\nu ^{(1)}(q\tau )(q\tau )^\nu ,$$ (37) where $`\tau `$ is the conformal time $$\tau _t^{\mathrm{}}\frac{dt^{}}{a(t^{})}=\frac{1}{a(t)H}.$$ (38) The contribution of the third kind to the expectation value then has the Fourier transform $`{\displaystyle d^{2\nu }xe^{i𝐪(𝐱𝐱^{})}\phi (𝐱,t)\phi (𝐱^{},t)_{\mathrm{𝐢𝐢𝐢}}}=\left({\displaystyle \frac{H^{2\nu 1}}{32\pi ^2}}\right)^3\left({\displaystyle \frac{2\pi }{q}}\right)^{4\nu }`$ $`\times 4\pi {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dp}{p}}{\displaystyle _{\mathrm{}}^t}𝑑t_1a^{2\nu }(t_1)\left({\displaystyle \frac{d}{dt_1}}\left\|(p\tau _1)^\nu H_\nu ^{(1)}(p\tau _1)\right\|^2\right)`$ $`\times \mathrm{Im}{\displaystyle \frac{d}{dt_1}}\left[\left((q\tau _1)^\nu H_\nu ^{(1)}(q\tau _1)\right)^2\left((q\tau )^\nu H_\nu ^{(1)}(q\tau )\right)^2\right]`$ (39) Let’s see what happens if we evaluate this by integrating first over $`p`$ and then over $`t_1`$ from $`\mathrm{}`$ to late times, or vice versa. To integrate first over $`p`$, we can change the variable of integration from $`p`$ to $`zp\tau _1`$, in which case the first derivative with respect to $`t_1`$ can be replaced with $`d/dt_1=(z/a_1\tau _1)(d/dz)=Hz(d/dz)`$, while $`dp/p=dz/z`$. Dimensional regularization (with $`2\nu <1`$) makes the function $`\left\|z^\nu H_\nu ^{(1)}(z)\right\|`$ vanish at $`z\mathrm{}`$, while for $`\nu >0`$ it takes the value $`2^\nu \mathrm{\Gamma }(\nu )/\pi `$ for $`z0`$, so $$_0^{\mathrm{}}𝑑z\frac{d}{dz}\left\|z^\nu H_\nu ^{(1)}(z)\right\|^2=\left(\frac{2^\nu \mathrm{\Gamma }(\nu )}{\pi }\right)^2,$$ and therefore $`{\displaystyle d^{2\nu }xe^{i𝐪(𝐱𝐱^{})}\phi (𝐱,t)\phi (𝐱^{},t)_{\mathrm{𝐢𝐢𝐢}}}=4\pi H\left({\displaystyle \frac{2^\nu \mathrm{\Gamma }(\nu )}{\pi }}\right)^2\left({\displaystyle \frac{H^{2\nu 1}}{32\pi ^2}}\right)^3\left({\displaystyle \frac{2\pi }{q}}\right)^{4\nu }`$ $`\times {\displaystyle _{\mathrm{}}^t}dt_1a^{2\nu }(t_1)\mathrm{Im}{\displaystyle \frac{d}{dt_1}}\left[\left((q\tau _1)^\nu H_\nu ^{(1)}(q\tau _1)\right)^2\left((q\tau )^\nu H_\nu ^{(1)}(q\tau )\right)^2\right]`$ For $`t_1+\mathrm{}`$ and $`t+\mathrm{}`$ (that is, $`\tau 0`$ and $`\tau _10`$), the integrand of the integral over $`t_1`$ on the second line has the constant limit $`a^{2\nu }(t_1)\mathrm{Im}{\displaystyle \frac{d}{dt_1}}\left[\left((q\tau _1)^\nu H_\nu ^{(1)}(q\tau _1)\right)^2\left((q\tau )^\nu H_\nu ^{(1)}(q\tau )\right)^2\right]{\displaystyle \frac{4\mathrm{\Gamma }(\nu )^2q^{2\nu }}{\pi ^3H^{2\nu 1}}}.`$ Thus for $`t\mathrm{}`$, the correlation function (39) does not approach a constant, but instead goes as $`{\displaystyle d^{2\nu }xe^{i𝐪(𝐱𝐱^{})}\phi (𝐱,t)\phi (𝐱^{},t)_{\mathrm{𝐢𝐢𝐢}}}{\displaystyle \frac{H^{4\nu 1}\mathrm{\Gamma }(\nu )^4t}{2(2\pi )^{104\nu }q^{2\nu }}}.`$ (42) There is no pole here that prevents continuation to space dimensionality $`2\nu =3`$. From this point of view, integrating first over $`p`$, the failure of the correlation function to approach a finite limit at late times is due to the fact already noted, that the integral over $`p`$ produces an effective interaction that does not satisfy the conditions of our theorem. But suppose we first integrate over $`t_1`$ from $`\mathrm{}`$ to $`+\mathrm{}`$. Now there is no problem with convergence at late times, because the original interaction does satisfy the conditions of our theorem, but instead we now have a problem with the convergence of the integral over $`p`$. It will be helpful to divide the integral over $`p`$ into an integral from $`0`$ to $`\mathrm{\Lambda }q`$, where $`\mathrm{\Lambda }1`$, and an integral from $`\mathrm{\Lambda }q`$ to infinity. The first integral obviously has no ultraviolet divergence, and the vanishing of the first time derivative in Eq. (39) for $`p0`$ prevents any infrared divergence. In the second integral $`p`$ and $`1/\tau `$ are the only magnitudes in the problem with which $`q`$ can be compared, so for $`t+\mathrm{}`$ and hence $`\tau 0`$ we can evaluate the correlation function by letting $`q0`$ and keeping only the leading term in $`q`$. Here again we can use the limiting formula (41), now for $`q0`$ instead of $`\tau 0`$ and $`\tau _10`$. The integral over $`t_1`$ is then trivial, and we find that for $`q1/\tau `$ the correlation function is $$d^{2\nu }xe^{i𝐪(𝐱𝐱^{})}\phi (𝐱,t)\phi (𝐱^{},t)_{\mathrm{𝐢𝐢𝐢}}\frac{H^{4\nu 2}\mathrm{\Gamma }(\nu )^4}{2(2\pi )^{102\nu }q^{2\nu }}_{\mathrm{\Lambda }q}^{\mathrm{}}\frac{dp}{p}+\mathrm{finite}.$$ (43) The ultraviolet divergent integral over $`p`$ is the price we pay for the naughtiness of taking the limit $`t\mathrm{}`$ before we integrate over $`p`$. In this model it is clear how to remedy the difficulty of calculating correlation functions at late times. As already mentioned, the original Lagrangian density (35) actually describes a free field theory. This is made manifest by defining a new scalar field $$\stackrel{~}{\phi }\sqrt{1+\lambda \phi ^2}𝑑\phi ,$$ (44) for which the Lagrangian density takes the form $$=\frac{1}{2}\sqrt{\mathrm{Detg}}g^{\mu \nu }_\mu \stackrel{~}{\phi }_\nu \stackrel{~}{\phi }.$$ (45) There is no problem in taking the late-time limit of the correlation function $`d^{2\nu }e^{i𝐪(𝐱𝐱^{})}\stackrel{~}{\phi }(𝐱,t)\stackrel{~}{\phi }(𝐱^{},t)`$ — it is just $`2^{2\nu }H^{2\nu 1}\mathrm{\Gamma }(\nu )^2/32\pi ^4q^{2\nu }`$. From this point of view, the growth of the correlation function (42) at late times is a result of our perversity in calculating the correlation function of $`\phi `$ instead of $`\stackrel{~}{\phi }`$. Can we find fields whose correlation functions have a constant limit at late times in theories that satisfy the conditions of our theorem but are not equivalent to free field theories? The general answer is not known, but here is a class of interacting field theories for which such “renormalized” fields can be found. This time we consider an arbitrary number of real scalar fields $`\phi _n(𝐱,t)`$ in a fixed de Sitter metric. The Lagrangian density is taken to have the form of a non-linear $`\sigma `$-model: $$=\frac{1}{2}\underset{nm}{}\sqrt{\mathrm{Detg}}g^{\mu \nu }\left(\delta _{nm}+\lambda K_{nm}(\phi )\right)_\mu \phi _n_\nu \phi _m,$$ (46) where $`K_{nm}(\phi )`$ is an arbitrary real symmetric matrix function of the $`\phi _n`$; $`\lambda `$ is a coupling constant; and again $`ae^{Ht}`$ with $`H`$ constant. The Hamiltonian derived from this Lagrangian density does satisfy the conditions of the theorem of Section IV, whatever the function $`K_{nm}(\phi )`$. To first order in $`\lambda `$, the same problem discussed earlier in this section arises from graphs in which an internal line of the field $`\phi _n`$ is emitted and absorbed from the same vertex, with a time derivative acting on just one end of this line. Depending on what correlation function is being calculated, the contribution of such graphs is proportional to various contractions of partial derivatives of the function $$A_m(\phi )\underset{n}{}\frac{K_{nm}(\phi )}{\phi _n}.$$ (47) Suppose we make a redefinition of the fields of first order in $`\lambda `$: $$\stackrel{~}{\phi }_n\phi _n\lambda \mathrm{\Delta }_n(\phi ).$$ (48) This changes the matrix $`K`$ to $$\stackrel{~}{K}_{nm}(\phi )=K_{nm}(\phi )+\frac{\mathrm{\Delta }_n(\phi )}{\phi _m}+\frac{\mathrm{\Delta }_m(\phi )}{\phi _n},$$ (49) and so $$\stackrel{~}{A}_m(\phi )\underset{n}{}\frac{\stackrel{~}{K}_{nm}(\phi )}{\phi _n}=A_m(\phi )+\underset{n}{}\frac{^2\mathrm{\Delta }_n(\phi )}{\phi _n\phi _m}+\underset{n}{}\frac{\mathrm{\Delta }_m(\phi )}{\phi _n\phi _n}.$$ (50) Thus the fields $`\stackrel{~}{\phi }_n`$ are renormalized, in the sense that to first order in $`\lambda `$ correlation functions have finite limits at late times, provided that $$\underset{n}{}\frac{^2\mathrm{\Delta }_n(\phi )}{\phi _n\phi _m}+\underset{n}{}\frac{\mathrm{\Delta }_m(\phi )}{\phi _n\phi _n}=A_m(\phi ).$$ (51) This can be solved by first solving the Poisson equation $$\underset{n}{}\frac{^2B(\phi )}{\phi _n\phi _n}=\frac{1}{2}\underset{n}{}\frac{A_n(\phi )}{\phi _n}$$ (52) and then solving a second Poisson equation $$\underset{n}{}\frac{^2\mathrm{\Delta }_m(\phi )}{\phi _n\phi _n}=A_m(\phi )\frac{B(\phi )}{\phi _m}.$$ (53) Thus for at least to first order in this class of theories, it is always possible to find a suitable set of renormalized fields. Because we can take the limit $`t\mathrm{}`$ only for the correlation functions of suitably defined fields (such as $`\stackrel{~}{\phi }_n`$ in our example), the question naturally arises, whether these are the fields whose correlation functions we want to calculate. The answer is conditioned by the fact that astronomical observations of the cosmic microwave background or large scale structure are made following a long era that has intervened since the end of inflation, during which things happened about which we know almost nothing, such as reheating, baryon and lepton synthesis, and dark matter decoupling. The only thing that allows us to use observations to learn about inflation is that some quantities were conserved during this era, while fluctuation wave lengths were outside the horizon. These are the only quantities whose correlation functions at the end of inflation can be interpreted in terms of current observations. In the classical limit, the quantities that are conserved outside the horizon are $`\zeta `$ and $`\gamma _{ij}`$, but we don’t know whether this will be true when quantum effects are taken into account. Still, we can expect that quantities are conserved only when there is some symmetry principle that makes them conserved, and whatever symmetry principle keeps some quantity conserved from the end of inflation to the time of horizon re-entry is likely also to keep it conserved from the time of horizon exit to the end of inflation. So we may guess that the quantities whose correlation functions we will need to know are just those whose correlation functions approach constant limits at the end of inflation. VI. DANGEROUS INTERACTIONS IN INFLATIONARY THEORIES We now return to the semi-realistic theories described in Section III. We will show in this section that all interactions are of the type called for in the theorem of Section IV; that is, they are all safe interactions that (in three space dimensions) do not grow exponentially at late times (and in fact are suppressed at late times at least by a factor $`a^1`$), or dangerous interactions containing only fields and not their time derivatives, which grow no faster that $`a`$ at late times. Fortunately, as noticed by Maldacena<sup>2</sup> in a different context, for this purpose it is not necessary to solve the constraint equations (7) and (8), which are quite complicated especially when the $`\sigma _n`$ fields are included. Inspection of these equations shows that when we count $`\dot{\zeta }`$, $`\dot{\gamma }_{ij}`$, and $`\dot{\sigma }_n`$ as of order $`a^2`$, the auxiliary fields $`N1`$ and $`N^i`$ are both also of order $`a^2`$.<sup>‡‡‡</sup><sup>‡‡‡</sup>‡‡‡In counting powers of $`a`$, note that the three-dimensional affine connection and Ricci tensor are independent of $`a`$, so the curvature scalar $`R^{(3)}`$ goes as $`a^2`$. For instance, for $`\gamma _{ij}=0`$, we have $`R^{(3)}=a^2e^{2\zeta }(4^2\zeta +2(\zeta )^2)`$. This is apparent in the first-order solution (9) of the constraint equations, but it holds to all orders in the fields. To calculate the quantity $`E^j{}_{i}{}^{}E_{}^{i}{}_{j}{}^{}(E^i{}_{i}{}^{})^2`$ in Eq. (5), we note that $$E^i{}_{j}{}^{}=H\delta ^i{}_{j}{}^{}+\dot{\zeta }\delta ^i{}_{j}{}^{}+\frac{1}{2}\left[e^\gamma \frac{}{t}e^\gamma \right]^i{}_{j}{}^{}\frac{1}{2}(^iN_j+_jN^i).$$ (54) The first term $`H\delta ^i_j`$ is of order zero in $`a`$, while all other terms are of order $`a^2`$, so $$E^j{}_{i}{}^{}E_{}^{i}{}_{j}{}^{}(E^i{}_{i}{}^{})^2=6H^212H\dot{\zeta }4H_kN^k+O(a^4)$$ (55) (In deriving this result, we note that $`\left[e^\gamma \frac{}{t}e^\gamma \right]^i{}_{i}{}^{}=\dot{\gamma }_{ii}=0`$.) The terms in (5) of first order in $`N1`$ all cancel as a consequence of the constraint equation (8), while terms of second order in $`N1`$ in Eq. (5) (and in particular in $`a^3e^{3\zeta }\dot{\overline{\phi }}^2/2N`$ and $`3H^3a^3e^{3\zeta }/2N`$) are suppressed by at least a factor $`a^3(a^2)^2`$, and are therefore safe. Therefore we can isolate all terms that are potentially dangerous by setting $`N=1`$, and find $`={\displaystyle \frac{a^3}{2}}e^{3\zeta }[R^{(3)}2V(\overline{\phi })6H^212H\dot{\zeta }4H_kN^k`$ $`+\dot{\overline{\phi }}^2a^2e^{2\zeta }[\mathrm{exp}(\gamma )]^{ij}{\displaystyle \underset{n}{}}_i\sigma _n_j\sigma _n]+O(a^1),`$ We note that $`e^{3\zeta }_kN^k=_k(e^{3\zeta }N^k)`$, so this term vanishes when integrated over three-space, and therefore makes no contribution to the action. The term proportional to $`\dot{\zeta }`$ can be written $$6a^3e^{3\zeta }H\dot{\zeta }=\frac{}{t}\left(2a^3He^{3\zeta }\right)+a^3e^{3\zeta }\left(6H^2+2\dot{H}\right).$$ The first term vanishes when integrated over time, so it gives no contribution to the action. To evaluate the remaining terms we use the unperturbed inflaton field equation, which (with $`8\pi G1`$) gives $`\dot{H}=\dot{\overline{\phi }}^2/2`$, and the Friedmann equation, which gives $`6H^2=2V+\dot{\overline{\phi }}^2`$. We then find a cancellation $$V3H^2+\frac{1}{2}\dot{\overline{\phi }}^2+6H^2+2\dot{H}=0.$$ Aside from terms that make no contribution to the action, the Lagrangian density is then $$=\frac{a^3}{2}e^{3\zeta }\left[R^{(3)}a^2e^{2\zeta }[\mathrm{exp}(\gamma )]^{ij}\underset{n}{}_i\sigma _n_j\sigma _n\right]+O(a^1).$$ (57) We see that, at least in this class of theories, the dangerous terms that are not suppressed by a factor $`a^1`$ grow at most like $`a`$ at late times, and involve only fields, not their time derivatives, as assumed in the theorem of Section III. It remains to be seen if in these theories, after integrating over times and taking the limit $`t\mathrm{}`$, the remaining integrals over internal wave numbers are made convergent by the same counterterms that eliminate ultraviolet divergences in flat spacetime, and if not, whether they can be made convergent by suitable redefinitions of the fields $`\zeta `$ and $`\gamma _{ij}`$ appearing in the correlation functions. This is left as a problem for further work. Not all theories satisfy the conditions of the theorem of Section IV. For instance, a non-derivative interaction $`\sqrt{\mathrm{Det}g}F(\sigma )`$ of the $`\sigma `$ fields would have $`+3`$ factors of $`a`$, and hence would violate the condition that the total number of factors of $`a`$ (counting each time derivative as -2 factors) must be no greater than $`+1`$. The $`\sigma `$ fields must be the Goldstone bosons of some broken global symmetry in order to satisfy the conditions of our theorem in a natural way. VII. A SAMPLE CALCULATION As an application of the formalism described in this paper, we will now calculate the one-loop contribution to the correlation function of two $`\zeta `$ fields, which is measured in the spectrum of anisotropies of the cosmic microwave background. As already mentioned, in the class of theories described in Section III, this two-point function is dominated by a matter loop, because there are many types of matter field and only one gravitational field. We first consider the contribution of second order in the interaction (27). It saves a great deal of work if we use the interaction-picture field equations (11) and (13) to put this interaction in the form $$H_{\zeta \sigma \sigma }(t)=d^3x_{\zeta \sigma \sigma }(𝐱,t)=A(t)+\dot{B}(t)$$ (58) where $`A`$ $`=`$ $`2ϵHa^5{\displaystyle \underset{n}{}}{\displaystyle d^3x\dot{\sigma }_n^2^2\dot{\zeta }}`$ (59) $`B`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle d^3x\left(\frac{a\zeta }{H}ϵa^3^2\dot{\zeta }\right)\left(\frac{1}{2}(\sigma _n)^2+\frac{1}{2}a^2\dot{\sigma }_n^2\right)}.`$ (60) In general, for an interaction Hamiltonian of the form (58), Eq. (2) can be put in the form $`Q(t)`$ $`=`$ $`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}i^N{\displaystyle _{\mathrm{}}^t}𝑑t_N{\displaystyle _{\mathrm{}}^{t_N}}𝑑t_{N1}\mathrm{}{\displaystyle _{\mathrm{}}^{t_2}}𝑑t_1`$ (61) $`\times [\stackrel{~}{H}_I(t_1),[\stackrel{~}{H}_I(t_2),\mathrm{}[\stackrel{~}{H}_I(t_N),\stackrel{~}{Q}^I(t)]\mathrm{}]],`$ where $$\stackrel{~}{H}_I(t)=e^{iB(t)}\left[A(t)+\dot{B}(t)+ie^{iB(t)}\left(\frac{d}{dt}e^{iB(t)}\right)\right]e^{iB(t)}=A(t)+i[B(t),A(t)]+\frac{i}{2}[B(t),\dot{B}(t)]+\mathrm{}$$ (62) $$\stackrel{~}{Q}^I(t)=e^{iB(t)}Q^I(t)e^{iB(t)}=Q^I(t)+i[B(t),Q^I(t)]\frac{1}{2}[B(t),[B(t),Q^I(t)]]+\mathrm{}.$$ (63) To second order in an interaction of the form (58), the expectation value is then $`Q(t)_2`$ $`=`$ $`{\displaystyle _{\mathrm{}}^t}𝑑t_2{\displaystyle _{\mathrm{}}^{t_2}}𝑑t_1[A(t_1),[A(t_2),Q^I(t)]]`$ (64) $`{\displaystyle _{\mathrm{}}^t}𝑑t_1[[B(t_1),A(t_1)+\dot{B}(t_1)/2],Q^I(t)]`$ $`[B(t),[B(t),Q^I(t)]],`$ The Fourier transform of the second-order term in the expectation value of a product of two $`\zeta `$s is then $`{\displaystyle d^3xe^{i𝐪(𝐱𝐱^{})}\mathrm{vac},\mathrm{in}\left\|\zeta (𝐱,t)\zeta (𝐱^{},t)\right\|\mathrm{vac},\mathrm{in}_2}`$ $`={\displaystyle \frac{32(2\pi )^9}{q^4}}\mathrm{Re}{\displaystyle _{\mathrm{}}^t}a^5(t_2)ϵ(t_2)H(t_2)𝑑t_2`$ $`\times {\displaystyle _{\mathrm{}}^{t_2}}a^5(t_1)ϵ(t_1)H(t_1)dt_1`$ $`\times \dot{\zeta }_q(t_1)\zeta _q^{}(t)\left(\dot{\zeta }_q(t_2)\zeta _q^{}(t)\zeta _q(t)\dot{\zeta }_q^{}(t_2)\right)`$ $`\times 𝒩{\displaystyle d^3pd^3p^{}\delta ^3(𝐩+𝐩^{}+𝐪)}`$ $`\times \dot{\sigma }_p(t_1)\dot{\sigma }_p^{}(t_2)\dot{\sigma }_p^{}(t_1)\dot{\sigma }_p^{}^{}(t_2)`$ $`+{\displaystyle \frac{(2\pi )^3}{4q^4}}𝒩{\displaystyle d^3pd^3p^{}\delta ^3(𝐩+𝐩^{}+𝐪)}`$ $`\times (𝐩𝐩^{})^2\|\sigma _p(t)\|^2\|\sigma _p^{}(t)\|^2`$ $`+\mathrm{}`$ (65) where $`𝒩`$ is the number of $`\sigma `$ fields. We have shown here explicitly the contribution of the first and third lines on the right-hand side of Eq. (64). The dots represent one-loop contributions of the second line, in which $`[B,A+\dot{B}/2]`$ plays the role of a $`\zeta \zeta \sigma \sigma `$ “seagull” interaction, as well as one-loop terms of first order in the $`\zeta \zeta \sigma \sigma `$ terms in Eq. (5), in both of which the integral over internal wave number is $`𝐪`$-independent, plus counterterms arising in first order from interactions that cancel ultraviolet divergences in flat space, including $`\sqrt{\mathrm{Det}g}R^{\mu \nu }R_{\mu \nu }`$ and $`\sqrt{\mathrm{Det}g}R^2`$ terms in the Lagrangian density that are not included in Eq. (5). Though it has not been made explicit in this section, we use dimensional regularization to remove infinities in the integrals over $`p`$ and $`p^{}`$ at intermediate stages in the calculation, and we now assume that the singularity as the number of space dimensions approaches three is cancelled by the terms in Eq. (65) represented by dots, leaving it to future work to show that this is the case. Then these integrals are dominated by $`pp^{}q`$. As we have seen, the integrals over time are then dominated by the time $`t_q`$ of horizon exit, when $`q/a(t_q)H(t_q)`$. For simplicity, we will assume (for the first time in this paper) that the unperturbed inflaton field $`\overline{\phi }(t)`$ is rolling very slowly down the potential at time $`t_q`$, so that the expansion near this time can be approximated as strictly exponential, $`a(t)e^{Ht}`$. Then the wave functions are $$\sigma _q(t)\sigma _q^oe^{iq\tau }\left(1+iq\tau \right),$$ $$\zeta _q(t)\zeta _q^oe^{iq\tau }\left(1+iq\tau \right),$$ where $`\tau `$ is the conformal time $$\tau _t^{\mathrm{}}\frac{dt}{a(t)},$$ and the wave functions outside the horizon have modulus $$\|\sigma _q^o\|^2=\frac{H^2(t_q)}{2(2\pi )^3q^3},\|\zeta _q^o\|^2=\frac{H^2(t_q)}{2(2\pi )^3ϵ(t_q)q^3}$$ Using these wave functions in Eq. (65) gives $`{\displaystyle d^3xe^{i𝐪(𝐱𝐱^{})}\mathrm{vac},\mathrm{in}\left\|\zeta (𝐱,t)\zeta (𝐱^{},t)\right\|\mathrm{vac},\mathrm{in}_2}`$ $`={\displaystyle \frac{(8\pi GH^2(t_q))^2𝒩}{(2\pi )^3}}{\displaystyle d^3pd^3p^{}\delta ^3(𝐩+𝐩^{}+𝐪)}`$ $`\times \left[{\displaystyle \frac{pp^{}}{q^7(p+p^{}+q)}}+{\displaystyle \frac{(𝐩𝐩^{})^2}{16q^4p^3p^3}}\right]+\mathrm{}`$ (66) with the dots having the same meaning as in Eq. (65). Simple dimensional analysis tells us that when the integral over internal wave numbers of the first term in square brackets is made finite by dimensional regularization, it is converted to $$d^3pd^3p^{}\delta ^3(𝐩+𝐩^{}+𝐪)\frac{pp^{}}{p+p^{}+q}q^{4+\delta }F(\delta ),$$ (67) where $`\delta `$ is a measure of the difference between the space dimensionality and three. The ultraviolet divergences in this integrals for $`\delta =0`$ gives the function $`F(\delta )`$ a singularities as $`\delta 0`$: $$F(\delta )\frac{F_0}{\delta }+F_1,$$ (68) so that in the limit $`\delta =0`$ $$d^3pd^3p^{}\delta ^3(𝐩+𝐩^{}+𝐪)\frac{pp^{}}{p+p^{}+q}=q^4\left[F_0\mathrm{ln}q+L\right],$$ (69) where $`L`$ is a divergent constant. We can easily calculate the coefficient $`F_0`$ of the logarithm. For this purpose, we note that, in general, $$d^3pd^3p^{}\delta ^3(𝐩+𝐩^{}+𝐪)f(p,p^{},q)=\frac{2\pi }{q}_0^{\mathrm{}}p𝑑p_{\|pq\|}^{p+q}p^{}𝑑p^{}f(p,p^{},q)$$ (70) To eliminate the divergence in the integral over $`p`$ and $`p^{}`$, we multiply by $`q`$ and differentiate six times with respect to $`q`$. A tedious but straightforward calculation gives $$\frac{d^6}{dq^6}\left[qd^3pd^3p^{}\delta ^3(𝐩+𝐩^{}+𝐪)\frac{pp^{}}{p+p^{}+q}\right]=\frac{8\pi }{q}$$ Comparing this with the result of applying the same operation to Eq. (69) then gives $`F_0=\pi /15`$. In contrast, the integral of the second term in square brackets in Eq. (66) is a sum of powers of $`q`$ with divergent coefficients, but with no logarithmic singularity in $`q`$. (This term would be eliminated if we calculated the expectation value of a product of fields $`\stackrel{~}{\zeta }\mathrm{exp}(iB)\zeta \mathrm{exp}(iB)`$ instead of $`\zeta `$.) The terms represented by dots in Eq. (65) make contributions that are also just a sum of powers of $`q`$ with divergent coefficients. We are assuming that all ultraviolet divergences cancel, but we cannot find resulting finite power terms without knowing the renormalized coefficients of the $`\sqrt{\mathrm{Det}g}R^{\mu \nu }R_{\mu \nu }`$ and $`\sqrt{\mathrm{Det}g}R^2`$ terms in the Lagrangian density. So we are left with the result (now restoring a suitable power of $`8\pi G`$) that $`{\displaystyle d^3xe^{i𝐪(𝐱𝐱^{})}\mathrm{vac},\mathrm{in}\left\|\zeta (𝐱,t)\zeta (𝐱^{},t)\right\|\mathrm{vac},\mathrm{in}_2}`$ $`={\displaystyle \frac{\pi \left(8\pi GH^2(t_q)\right)^2𝒩}{15(2\pi )^3q^3}}\left[\mathrm{ln}q+C\right]`$ (71) with $`C`$ an unknown constant. This may be compared with the classical (and classic) result, that in slow roll inflation this correlation function takes the form $`{\displaystyle d^3xe^{i𝐪(𝐱𝐱^{})}\mathrm{vac},\mathrm{in}\left\|\zeta (𝐱,t)\zeta (𝐱^{},t)\right\|\mathrm{vac},\mathrm{in}_0}`$ $`={\displaystyle \frac{8\pi GH^2(t_q)}{4(2\pi )^3\|ϵ(t_q)\|q^3}}`$ (72) The one-loop correction (71) is smaller by a factor of order $`8\pi GH^2𝒩\|ϵ(t_q)\|`$, so even if $`𝒩`$ is $`10^2`$ or $`10^3`$ this correction is likely to remain unobservable. Still, it is interesting that even in the extreme slow roll limit, where $`H(t_q)`$ and $`ϵ(t_q)`$ are nearly constant, the factor $`\mathrm{ln}q`$ gives it a different dependence on the wave number $`q`$. ACKNOWLEDGMENTS For helpful conversations I am grateful to K. Chaicherdsakul, S. Deser, W. Fischler, E. Komatsu, J. Maldacena, A. Vilenkin, and R. Woodard. This material is based upon work supported by the National Science Foundation under Grants Nos. PHY-0071512 and PHY-0455649 and with support from The Robert A. Welch Foundation, Grant No. F-0014, and also grant support from the US Navy, Office of Naval Research, Grant Nos. N00014-03-1-0639 and N00014-04-1-0336, Quantum Optics Initiative. APPENDIX: THE IN-IN FORMALISM 1. Time Dependence First, it is necessary to be precise about the origin of the time-dependence of the fluctuation Hamiltonian in applications such as those encountered in cosmology. Consider a general Hamiltonian system, with canonical variables $`\varphi _a(𝐱,t)`$ and conjugates $`\pi _a(𝐱,t)`$ satisfying the commutation relations $$[\varphi _a(𝐱,t),\pi _b(𝐲,t)]=i\delta _{ab}\delta ^3(𝐱𝐲),[\varphi _a(𝐱,t),\varphi _b(𝐲,t)]=[\pi _a(𝐱,t),\pi _b(𝐲,t)]=0,$$ (A.1) and the equations of motion $$\dot{\varphi }_a(𝐱,t)=i[H[\varphi (t),\pi (t)],\varphi _a(𝐱,t)],\dot{\pi }_a(𝐱,t)=i[H[\varphi (t),\pi (t)],\pi _a(𝐱,t)].$$ (A.2) Here $`a`$ is a compound index labeling particular fields and their spin components. The Hamiltonian $`H`$ is a functional of the $`\varphi _a(𝐱,t)`$ and $`\pi _a(𝐱,t)`$ at fixed time $`t`$, which according to Eq. (A.2) is of course independent of the time at which these variables are evaluated. We assume the existence of a time-dependent c-number solution $`\overline{\varphi }_a(𝐱,t)`$, $`\overline{\pi }_a(𝐱,t)`$, satisfying the classical equations of motion: $$\dot{\overline{\varphi }}_a(𝐱,t)=\frac{\delta H[\overline{\varphi }(t),\overline{\pi }(t)]}{\delta \overline{\pi }_a(𝐱,t)},\dot{\overline{\pi }}_a(𝐱,t)=\frac{\delta H(\overline{\varphi }(t),\overline{\pi }(t)]}{\delta \overline{\varphi }_a(𝐱,t)},$$ (A.3) and we expand around this solution, writing $$\varphi _a(𝐱,t)=\overline{\varphi }_a(𝐱,t)+\delta \varphi _a(𝐱,t),\pi _a(𝐱,t)=\overline{\pi }_a(𝐱,t)+\delta \pi _a(𝐱,t).$$ (A.4) (In cosmology, $`\overline{\varphi }_a`$ would describe the Robertson–Walker metric and the expectation values of various scalar fields.) Of course, since c-numbers commute with everything, the fluctuations satisfy the same commutation rules (A.1) as the total variables: $$[\delta \varphi _a(𝐱,t),\delta \pi _b(𝐲,t)]=i\delta _{ab}\delta ^3(𝐱𝐲),[\delta \varphi _a(𝐱,t),\delta \varphi _b(𝐱,t)]=[\delta \pi _a(𝐱,t),\delta \pi _b(𝐱,t)]=0,$$ (A.5) When the Hamiltonian is expanded in powers of the perturbations $`\delta \varphi _a(𝐱,t)`$ and $`\delta \pi _a(𝐱,t)`$ at some definite time $`t`$, we encounter terms of zeroth and first order in the perturbations, as well as time-dependent terms of second and higher order: $`H[\varphi (t),\pi (t)]`$ $`=`$ $`H[\overline{\varphi }(t),\overline{\pi }(t)]+{\displaystyle \underset{a}{}}{\displaystyle \frac{\delta H[\overline{\varphi }(t),\overline{\pi }(t)]}{\delta \overline{\varphi }_a(𝐱,t)}}\delta \varphi _a(𝐱,t]+{\displaystyle \underset{a}{}}{\displaystyle \frac{\delta H[\overline{\varphi }(t),\overline{\pi }(t)]}{\overline{\pi }_a(𝐱,t)}}\delta \pi _a(𝐱,t)`$ (A.6) $`+\stackrel{~}{H}[\delta \varphi (t),\delta \pi (t);t],`$ where $`\stackrel{~}{H}[\delta \varphi (t),\delta \pi (t);t]`$ is the sum of all terms in $`H[\overline{\varphi }(t)+\delta \varphi (t),\overline{\pi }(t)+\delta \pi (t)]`$ of second and higher order in the $`\delta \varphi (𝐱,t)`$ and/or $`\delta \pi (𝐱,t)`$. Now, although $`H`$ generates the time-dependence of $`\varphi _a(𝐱,t)`$ and $`\pi _a(𝐱,t)`$, it is $`\stackrel{~}{H}`$ rather than $`H`$ that generates the time dependence of $`\delta \varphi _a(𝐱,t)`$ and $`\delta \pi _a(𝐱,t)`$. That is, Eq. (A.2) gives $$\dot{\overline{\varphi }}_a(𝐱,t)+\delta \dot{\varphi }_a(𝐱,t)=i[H[\varphi (t),\pi (t)],\delta \varphi _a(𝐱,t)],\dot{\overline{\pi }}_a(𝐱,t)+\delta \dot{\pi }_a(𝐱,t)=i[H[\varphi (t),\pi (t)],\delta \pi _a(𝐱,t)],$$ while Eqs. (A.5) and (A.3) give $$i[\underset{b}{}d^3y\frac{\delta H[\overline{\varphi }(t),\overline{\pi }(t)]}{\delta \overline{\varphi }_b(𝐲,t)}\delta \varphi _b(𝐲,t)+\underset{b}{}d^3y\frac{\delta H[\overline{\varphi }(t),\overline{\pi }(t)]}{\delta \overline{\pi }_b(𝐲,t)}\delta \pi _b(𝐲,t),\delta \varphi _a(𝐱,t)]=\dot{\overline{\varphi }}_a(𝐱,t)$$ $$i[\underset{b}{}d^3y\frac{\delta H[\overline{\varphi }(t),\overline{\pi }(t)]}{\delta \overline{\varphi }_b(𝐲,t)}\delta \varphi _b(𝐲,t)+\underset{b}{}d^3y\frac{\delta H[\overline{\varphi }(t),\overline{\pi }(t)]}{\delta \overline{\pi }_b(𝐲,t)}\delta \pi _b(𝐲,t),\delta \pi _a(𝐱,t)]=\dot{\overline{\pi }}_a(𝐱,t).$$ Subtracting, we find $$\delta \dot{\varphi }_a(𝐱,t)=i[\stackrel{~}{H}[\varphi (t),\pi (t);t],\delta \varphi _a(𝐱,t)],\delta \dot{\pi }_a(𝐱,t)=i[\stackrel{~}{H}[\varphi (t),\pi (t);t],\delta \pi _a(𝐱,t)].$$ (A.7) This then is our prescription for constructing the time-dependent Hamiltonian $`\stackrel{~}{H}`$ that governs the time-dependence of the fluctuations: expand the original Hamiltonian $`H`$ in powers of fluctuations $`\delta \varphi `$ and $`\delta \pi `$, and throw away the terms of zeroth and first order in these fluctuations. It is this construction that gives $`\stackrel{~}{H}`$ an explicit dependence on time. 2. Operator Formalism for Expectation Values We consider a general Hamiltonian system, of the sort described in the previous subsection. It follows from Eq. (A.7) that the fluctuations at time $`t`$ can be expressed in terms of the same operators at some very early time $`t_0`$ through a unitary transformation $$\delta \varphi _a(t)=U^1(t,t_0)\delta \varphi _a(t_0)U(t,t_0),\delta \pi _a(t)=U^1(t,t_0)\delta \pi _a(t_0)U(t,t_0),$$ (A.8) where $`U(t,t_0)`$ is defined by the differential equation $$\frac{d}{dt}U(t,t_0)=i\stackrel{~}{H}[\delta \varphi (t),\delta \pi (t);t]U(t,t_0)$$ (A.9) and the initial condition $$U(t_0,t_0)=1.$$ (A.10) In the application that concerns us in cosmology, we can take $`t_0=\mathrm{}`$, by which we mean any time early enough so that the wavelengths of interest are deep inside the horizon. To calculate $`U(t,t_0)`$, we now further decompose $`\stackrel{~}{H}`$ into a kinematic term $`H_0`$ that is quadratic in the fluctuations, and an interaction term $`H_I`$: $$\stackrel{~}{H}[\delta \varphi (t),\delta \pi (t);t]=H_0[\delta \varphi (t),\delta \pi (t);t]+H_I[\delta \varphi (t),\delta \pi (t);t],$$ (A.11) and we seek to calculate $`U`$ as a power series in $`H_I`$. To this end, we introduce an “interaction picture”: we define fluctuation operators $`\delta \varphi _a^I(t)`$ and $`\delta \pi _a^I(t)`$ whose time dependence is generated by the quadratic part of the Hamiltonian: $$\delta \dot{\varphi }_a^I(t)=i[H_0[\delta \varphi ^I(t),\delta \pi ^I(t);t],\delta \varphi _a^I(t)],\delta \dot{\pi }_a^I(t)=i[H_0[\delta \varphi ^I(t),\delta \pi ^I(t);t],\delta \pi _a^I(t)],$$ (A.12) and the initial conditions $$\delta \varphi _a^I(t_0)=\delta \varphi _a(t_0),\delta \pi _a^I(t_0)=\delta \pi _a(t_0).$$ (A.13) Because $`H_0`$ is quadratic, the interaction picture operators are free fields, satisfying linear wave equations. It follows from Eq. (A.12) that in evaluating $`H_0[\delta \varphi ^I,\delta \pi ^I;t]`$ we can take the time argument of $`\delta \varphi ^I`$ and $`\delta \pi ^I`$ to have any value, and in particular we can take it as $`t_0`$, so that $$H_0[\delta \varphi ^I(t),\delta \pi ^I(t);t]=H_0[\delta \varphi (t_0),\delta \pi (t_0);t],$$ (A.14) but the intrinsic time-dependence of $`H_0`$ still remains. The solution of Eq. (A.12) can again be written as a unitary transformation: $$\delta \varphi _a^I(t)=U_0^1(t,t_0)\delta \varphi _a(t_0)U_0(t,t_0),\delta \pi _a^I(t)=U_0^1(t,t_0)\delta \pi _a(t_0)U_0(t,t_0),$$ (A.15) with $`U_0`$ defined by the differential equation $$\frac{d}{dt}U_0(t,t_0)=iH_0[\delta \varphi (t_0),\delta \pi (t_0);t]U_0(t,t_0)$$ (A.16) and the initial condition $$U_0(t_0,t_0)=1.$$ (A.17) Then from Eqs. (A.9) and (A.16) we have $$\frac{d}{dt}\left[U_0^1(t,t_0)U(t,t_0)\right]=iU_0^1(t,t_0)H_I[\delta \varphi (t_0),\delta \pi (t_0);t]U(t,t_0).$$ Using Eq. (A.15), this gives $$U(t,t_0)=U_0(t,t_0)F(t,t_0),$$ (A.18) where $$\frac{d}{dt}F(t,t_0)=iH_I(t)F(t,t_0),F(t_0,t_0)=1.$$ (A.19) and $`H_I(t)`$ is the interaction Hamiltonian in the interaction picture: $$H_I(t)U_0(t,t_0)H_I[\delta \varphi (t_0),\delta \pi (t_0);t]U_0^1(t,t_0)=H_I[\delta \varphi ^I(t),\delta \pi ^I(t);t]$$ (A.20) The solution of equations like (A.19) is well known $$F(t,t_0)=T\mathrm{exp}\left(i_{t_0}^tH_I(t)𝑑t\right)$$ (A.21) where $`T`$ indicates that the products of $`H_I`$s in the power series expansion of the exponential are to be time-ordered; that is, they are to be written from left to right in the decreasing order of time arguments. The solution for the fluctuations in terms of the free fields of the interaction picture is then given by Eqs. (A.8) and (A.15) as $`Q(t)`$ $`=`$ $`F^1(t,t_0)Q^I(t)F(t,t_0)`$ (A.22) $`=`$ $`\left[\overline{T}\mathrm{exp}\left(i{\displaystyle _{t_0}^t}H_I(t)𝑑t\right)\right]Q^I(t)\left[T\mathrm{exp}\left(i{\displaystyle _{t_0}^t}H_I(t)𝑑t\right)\right],`$ where $`Q(t)`$ is any $`\delta \varphi (𝐱,t)`$ or $`\delta \pi (𝐱,t)`$ or any product of the $`\delta \varphi `$s and/or $`\delta \pi `$s, all at the same time $`t`$ but in general with different space coordinates, and $`Q^I(t)`$ is the same product of $`\delta \varphi ^I(𝐱,t)`$ and/or $`\delta \pi ^I(𝐱,t)`$. Also, $`\overline{T}`$ denotes anti-time-ordering: products of $`H_I`$s in the power series expansion of the exponential are to be written from left to right in the increasing order of time arguments. 3. Diagrammatic Formalism for Expectation Values We want to use Eq. (A.22) to calculate the expectation value $`Q(t)`$of the product $`Q(t)`$ in a “Bunch–Davies” vacuum, annihilated by the positive-frequency part of the interaction picture fluctuations $`\delta \phi ^I`$ and $`\delta \pi ^I`$. We can use the familiar Wick theorem to express the vacuum expectation value of the right-hand side of Eq. (A.22) as a sum over pairings of the $`\delta \phi ^I`$ and $`\delta \pi ^I`$ with each other. (This of course is the same as supposing the interaction-picture fields in $`H_I(t)`$ and $`Q^I(t)`$ to be governed by a Gaussian probability distribution, except that the order of operators in bilinear averages has to be the same as the order in which they appear in Eq. (A.22).) Expanding Eq. (A.22) as a sum of products of bilinear products leads to a set of diagrammatic rules, but one that is rather complicated. In calculating the term in $`Q`$ of $`N`$th order in the interaction, we draw all diagrams with $`N`$ vertices. Just as for ordinary Feynman diagrams, each vertex is labeled with a space and time coordinate, and has lines attached corresponding to the fields in the interaction. There are also external lines, one for each field operator in the product $`Q`$, labeled with the different space coordinates and the common time $`t`$ in the arguments of these fields. All external lines are connected to vertices or other external lines, and all remaining lines attached to vertices are attached to other vertices. But there are significant differences between the rules following from Eq. (A.22) and the usual Feynman rules: * We have to distinguish between “right” and “left” vertices, arising respectively from the time-ordered product and the anti-time-ordered product. A diagram with $`N`$ vertices contributes a sum over all $`2^N`$ ways of choosing each vertex to be a left vertex or a right vertex. Each right or left vertex contributes a factor $`i`$ or $`+i`$, respectively, as well as whatever coupling parameters appear in the interaction. * A line connecting two right vertices or a right vertex and an external line, in which it is associated with field operators $`A(𝐱,t^{})`$ and $`B(𝐲,t^{\prime \prime })`$, contributes a conventional Feynman propagator $`T\{A(𝐱,t^{})B(𝐲,t^{\prime \prime }\}`$. (It will be understood here and below, that in calculating propagators all fields $`A`$, $`B`$, etc. are taken in the interaction picture, and can be $`\delta \phi ^I`$s and/or $`\delta \pi ^I`$s.) As a special case, if $`B`$ is associated with an external line then $`t^{\prime \prime }=t`$, and since $`t^{}t`$, this is $`B(𝐲,t)A(𝐱,t^{})`$. * A line connecting two left vertices, associated with field operators $`A(𝐱,t^{})`$ and $`B(𝐲,t^{\prime \prime })`$, contributes a propagator $`\overline{T}\{A(𝐱,t^{})B(𝐲,t^{\prime \prime }\}`$. As a special case, if $`B`$ is associated with an external line then $`t^{\prime \prime }=t`$, and this is $`A(𝐱,t^{})B(𝐲,t)`$. * A line connecting a left vertex, in which it is associated with a field operator $`A(𝐱,t^{})`$, to a right vertex, in which it is associated with a field operator $`B(𝐲,t^{\prime \prime })`$, contributes a propagator $`A(𝐱,t^{})B(𝐲,t^{\prime \prime })`$. * We must integrate over all over the times $`t^{},t^{\prime \prime },\mathrm{}`$, associated with the vertices from $`t_0`$ to $`t`$, as well as over all space coordinates associated with the vertices. We must say a word about the disconnected parts of diagrams. A vacuum fluctuation subdiagram is one in which each vertex is connected only to other vertices, not to external lines. Just as in ordinary quantum field theories, the sum of all vacuum fluctuation diagrams contributes a numerical factor multiplying the contribution of diagrams in which vacuum fluctuations are excluded. But unlike the case of ordinary quantum field theory, this numerical factor is not a phase factor, but is simply $$\left[\overline{T}\mathrm{exp}\left(i_{t_0}^tH_I(t)𝑑t\right)\right]\left[T\mathrm{exp}\left(i_{t_0}^tH_I(t)𝑑t\right)\right]=1.$$ (A.23) Hence in the “in-in” formalism all vacuum fluctuation diagrams automatically cancel. Even so, a diagram may contain disconnected parts which do not cancel, such as external lines passing through the diagram without interacting. Ignoring all disconnected parts gives what in the theory of noise is known as the cumulants of expectation values,<sup>10</sup> from which the full expectation values can easily be calculated as a sum of products of cumulants. 4. Path Integral Derivation of the Diagrammatic Rules. It is often preferable use path integration instead of the operator formalism, in order to derive the Feynman rules directly from the Lagrangian rather than from the Hamltonian, or to make available a larger range of gauge choices, or to go beyond perturbation theory. Going back to Eq. (1), and following the same reasoning<sup>11</sup> that leads from the operator formalism to the path-integral formalism in the calculation of S-matrix elements, we see that the vacuum expectation value of any product $`Q(t)`$ of $`\delta \varphi `$s and $`\delta \pi `$s at the same time $`t`$ (now taking $`t_0=\mathrm{}`$) is $`Q(t)`$ $`=`$ $`{\displaystyle \underset{𝐱,t^{},a}{}d\delta \varphi _{La}(𝐱,t^{})\underset{𝐱,t^{},a}{}\frac{d\delta \pi _{La}(𝐱,t^{})}{2\pi }\underset{𝐱,t^{},a}{}d\delta \varphi _{Ra}(𝐱,t^{})\underset{𝐱,t^{},a}{}\frac{d\delta \pi _{Ra}(𝐱,t^{})}{2\pi }}`$ (A.24) $`\times \mathrm{exp}\left\{i{\displaystyle _{\mathrm{}}^t}𝑑t^{}\left[{\displaystyle \underset{a}{}}{\displaystyle d^3x\delta \dot{\varphi }_{La}(𝐱,t^{})\delta \pi _{La}(𝐱,t^{})}\stackrel{~}{H}[\delta \varphi _L(t^{}),\delta \pi _L(t^{});t^{}]\right]\right\}`$ $`\times \mathrm{exp}\left\{i{\displaystyle _{\mathrm{}}^t}𝑑t^{}\left[{\displaystyle \underset{a}{}}{\displaystyle d^3x\delta \dot{\varphi }_{Ra}(𝐱,t^{})\delta \pi _{Ra}(𝐱,t^{})}\stackrel{~}{H}(\delta \varphi _R(𝐱,t^{}),\delta \pi _R(𝐱,t^{});t^{})\right]\right\}`$ $`\times {\displaystyle \underset{𝐱,a}{}}\delta (\delta \varphi _{La}(𝐱,t)\delta \varphi _{Ra}(𝐱,t))\delta (\delta \pi _{La}(𝐱,t)\delta \pi _{Ra}(𝐱,t))Q[\delta \varphi _L(t),\delta \pi _L(t)]`$ $`\times \mathrm{\Psi }_0^{}\left[\delta \varphi _L(\mathrm{})\right]\mathrm{\Psi }_0\left[\delta \varphi _R(\mathrm{})\right].`$ Here the functional $`\mathrm{\Psi }_0[\delta \varphi ]`$ is the wave function of the vacuum,<sup>12</sup> $`\mathrm{\Psi }_0[\varphi (\mathrm{})]\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{a,b}{}}{\displaystyle d^3xd^3y_{ab}(𝐱,𝐲)\delta \varphi _a(𝐱,\mathrm{})\delta \varphi _b(𝐲,\mathrm{})}\right)`$ $`=\mathrm{exp}\left({\displaystyle \frac{ϵ}{2}}{\displaystyle _{\mathrm{}}^t}𝑑t^{}e^{ϵt^{}}{\displaystyle \underset{a,b}{}}{\displaystyle d^3xd^3y_{ab}(𝐱,𝐲)\delta \varphi _a(t^{})\delta \varphi _b(t^{})}\right),`$ (A.25) where $`_{ab}`$ is a positive-definite kernel. For instance, for a real scalar field of mass $`m`$, $$(𝐱,𝐲)\frac{1}{(2\pi )^3}d^3pe^{i𝐩(𝐱𝐲)}\sqrt{𝐩^2+m^2}.$$ (A.26) As is well known, if the Hamiltonian is quadratic in the canonical conjugates $`\delta \pi _a`$ with a field-independent coefficient in the term of second order, then we can integrate over the $`\delta \pi _a`$ by simply setting $`\delta \dot{\varphi }_a=\stackrel{~}{H}/\delta \pi _a`$, and the quantity $`_a\delta \dot{\varphi }_a(t^{})\delta \pi _a(t^{})\stackrel{~}{H}(\delta \varphi (t^{}),\delta \pi (t^{});t^{})`$ in Eq. (A.24) then becomes the original Lagrangian. We will not pursue this here, but will rather take up a puzzle that at first sight seems to throw doubt on the equivalence of the path integral formula (A.24), when we do not integrate out the $`\pi `$s, with the operator formalism. The puzzle is that, although the propagators for lines connecting left vertices to each other or right vertices to each other or left or right vertices to external lines are Greens functions of the sort that familiarly emerge from path integrals, what are we to make of the propagators arising from Eq. (A.22) for lines connecting left vertices with right vertices? These are not Greens functions; that is, they are solutions of homogeneous wave equations, not of inhomogeneous wave equations with a delta function source. As we shall see, the source of these propagators lies in the delta functions in Eq. (A.24). It is these delta functions that tie together the integrals over the $`L`$ variables and over the $`R`$ variables, so that the expression (A.18) does not factor into a product of these integrals. In analyzing the consequences of Eq. (A.24), it is convenient to condense our notation yet further, and let a variable $`\xi _n(t)`$ stand for all the $`\delta \varphi _a(𝐱,t)`$ and $`\delta \pi _a(𝐱,t)`$, so that $`n`$ runs over positions in space and whatever discrete indices are used to distinguish different fields, plus a two-valued index that distinguishes $`\delta \varphi `$ from $`\delta \pi `$. With this understanding, Eq. (A.24) reads $`Q(t)={\displaystyle \underset{t^{},n}{}\frac{d\xi _{Ln}(t^{})}{\sqrt{2\pi }}\underset{t^{},n}{}\frac{d\xi _{Rn}(t^{})}{\sqrt{2\pi }}}`$ $`\times \mathrm{exp}\left\{i{\displaystyle _{\mathrm{}}^t}𝑑t^{}\stackrel{~}{L}(\xi _L(t^{}),\dot{\xi }_L(t^{});t^{})\right\}\mathrm{exp}\left\{i{\displaystyle _{\mathrm{}}^t}𝑑t^{}\stackrel{~}{L}[\xi _R(t^{}),\dot{\xi }_R(t^{});t^{}]\right\}`$ $`\times \left({\displaystyle \underset{n}{}}\delta \left(\xi _{Ln}(t)\xi _{Rn}(t)\right)\right)Q\left(\xi _L(t)\right)\mathrm{\Psi }_0^{}\left(\xi _L(\mathrm{})\right)\mathrm{\Psi }_0\left(\xi _R(\mathrm{})\right),`$ (A.27) where $$\stackrel{~}{L}[\xi (t^{}),\dot{\xi }(t^{});t^{}]\underset{a}{}d^3x\delta \pi _a(𝐱,t^{})\delta \dot{\varphi }_a(𝐱,t^{})\stackrel{~}{H}[\delta \varphi (t^{}),\delta \pi (t^{});t^{}].$$ (A.28) To expand in powers of the interaction, we split $`\stackrel{~}{L}`$ into a term $`\stackrel{~}{L}_0`$ that is quadratic in the fluctuations, plus an interaction term $`\stackrel{~}{H}_I`$: $$\stackrel{~}{L}=\stackrel{~}{L}_0\stackrel{~}{H}_I,$$ (A.29) where $$\stackrel{~}{L}_0[\xi (t^{}),\dot{\xi }(t^{});t^{}]=\underset{a}{}d^3x\delta \dot{\varphi }_a(𝐱,t^{})\delta \pi _a(𝐱,t^{})\stackrel{~}{H}_0(\delta \varphi (t^{}),\delta \pi (t^{});t^{}).$$ (A.30) As in calculations of the S-matrix, we will include the argument of the exponential in the vacuum wave functions along with the quadratic part of the Lagrangian, writing $`{\displaystyle _{\mathrm{}}^t}dt^{}\{\stackrel{~}{L}_0[\xi _R(t^{}),\dot{\xi }_R(t^{});t^{}]`$ $`+{\displaystyle \frac{iϵ}{2}}{\displaystyle \underset{ab}{}}{\displaystyle }d^3x{\displaystyle }d^3y_{ab}(𝐱,𝐲)\delta \varphi _{Ra}(𝐱,t^{})\delta \varphi _{Rb}(𝐲,t^{})\}`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{nn^{}}{}}{\displaystyle \underset{t^{},t^{\prime \prime }}{}}𝒟_{nt^{},mt^{\prime \prime }}^R\xi _{Rn}(t^{})\xi _{Rn^{}}(t^{\prime \prime }),`$ (A.31) $`{\displaystyle _{\mathrm{}}^t}dt^{}\{\stackrel{~}{L}_0[\xi _L(t^{}),\dot{\xi }_L(t^{});t^{}]`$ $`{\displaystyle \frac{iϵ}{2}}{\displaystyle \underset{ab}{}}{\displaystyle }d^3x{\displaystyle }d^3y_{ab}(𝐱,𝐲)\delta \varphi _{La}(𝐱,t^{})\delta \varphi _{Lb}(𝐲,t^{})\}`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{nn^{}}{}}{\displaystyle \underset{t^{},t^{\prime \prime }}{}}𝒟_{nt^{},n^{}t^{\prime \prime }}^L\xi _{Ln}(t^{})\xi _{Ln^{}}(t^{\prime \prime })`$ (A.32) The vacuum wave function is the same for $`\xi _L`$ and $`\xi _R`$, but it is combined here with an exponential $`\mathrm{exp}(i\stackrel{~}{L}_0)`$ for the $`\xi _{Ln}`$ and an exponential $`\mathrm{exp}(+i\stackrel{~}{L}_0)`$ for the $`\xi _{Rn}`$, which accounts for the different signs of the $`iϵ`$ terms in Eqs. (A.31) and (A.32). (The factor $`e^{ϵt^{}}`$ in Eq. (A.25) is effectively equal to one for any finite $`t^{}`$, and has therefore been dropped.) We also express the product of delta functions in Eq. (A.27) as a Gaussian: $`{\displaystyle \underset{n}{}}\delta \left(\xi _{Ln}(t)\xi _{Rn}(t)\right)\mathrm{exp}\left({\displaystyle \frac{1}{ϵ^{}}}{\displaystyle \underset{n}{}}\left(\xi _{Ln}(t)\xi _{Rn}(t)\right)^2\right)`$ $`=\mathrm{exp}\left({\displaystyle \underset{nn^{}}{}}{\displaystyle \underset{t^{}t^{\prime \prime }}{}}𝒞_{nt^{},n^{}t^{\prime \prime }}\left(\xi _{Ln}(t^{})\xi _{Rn}(t^{})\right)\left(\xi _{Ln^{}}(t^{\prime \prime })\xi _{Rn^{}}(t^{\prime \prime })\right)\right),`$ (A.33) where $$𝒞_{nt^{},n^{}t^{\prime \prime }}\frac{1}{ϵ^{}}\delta _{nn^{}}\delta (t^{}t)\delta (t^{\prime \prime }t),$$ (A.34) and $`ϵ^{}`$ is another positive infinitesimal. Following the usual rules for integrating a Gaussian times a polynomial, the integral is given by a sum over diagrams as described above, but with a line that connects right vertices with each other (or with external lines) contributing a factor $`i\mathrm{\Delta }_{nt^{},n^{}t^{\prime \prime }}^{RR}`$, a line that connects left vertices with each other (or with external lines) contributing a factor $`i\mathrm{\Delta }_{nt^{},n^{}t^{\prime \prime }}^{LL}`$, and a line that connects a right vertex where it is associated with $`\xi _n(t^{})`$ with a left vertex associated with $`\xi _n^{}(t^{\prime \prime })`$ contributing a factor $`i\mathrm{\Delta }_{nt^{},n^{}t^{\prime \prime }}^{RL}`$, with the $`\mathrm{\Delta }`$s determined by the condition $$\left(\begin{array}{cc}i𝒟^R𝒞& 𝒞\\ 𝒞& i𝒟^L𝒞\end{array}\right)\left(\begin{array}{cc}i\mathrm{\Delta }^{RR}& i\mathrm{\Delta }^{RL}\\ i(\mathrm{\Delta }^{RL})^\mathrm{T}& i\mathrm{\Delta }^{LL}\end{array}\right)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$ (A.35) This must hold whatever tiny value we give to $`ϵ^{}`$, and so $$𝒟^R\mathrm{\Delta }^{RR}=1,𝒟^L\mathrm{\Delta }^{LL}=1,$$ (A.36) $$𝒟^R\mathrm{\Delta }^{RL}=0,𝒟^L\left(\mathrm{\Delta }^{RL}\right)^\mathrm{T}=0,$$ (A.37) $$𝒞\mathrm{\Delta }^{LL}=𝒞\mathrm{\Delta }^{RL},𝒞\mathrm{\Delta }^{RR}=𝒞(\mathrm{\Delta }^{RL})^\mathrm{T}.$$ (A.38) The first Eq. (A.36) is the usual inhomogeneous wave equation for the propagator, whose solution as well known is $$i\mathrm{\Delta }_{nt^{},n^{}t^{\prime \prime }}^{RR}=T\{\xi _n(t^{})\xi _n^{}(t^{\prime \prime })\},$$ (A.39) with the time-ordering dictated by the $`+iϵ`$ in Eq. (A.31). The second Eq. (A.36) is the complex conjugate of the first wave equation, whose solution is the complex conjugate of Eq. (A.39): $$i\mathrm{\Delta }_{nt^{},n^{}t^{\prime \prime }}^{LL}=\overline{T}\{\xi _n(t^{})\xi _n^{}(t^{\prime \prime })\}.$$ (A.40) Eqs. (A.39) and (A.40) thus give the same propagators for lines connecting right vertices with each other or with external lines, and for lines connecting left vertices with each other or with external lines, as we we encountered in the operator formalism. Equations (A.37) tell us that $`\mathrm{\Delta }^{RL}`$ and $`(\mathrm{\Delta }^{RL})^\mathrm{T}`$ satisfy the homogeneous versions of the wave equations satisfied by $`\mathrm{\Delta }^{RR}`$ and $`\mathrm{\Delta }^{LL}`$, but to find $`\mathrm{\Delta }^{RL}`$ we also need an initial condition. This is provided by the first of Eqs. (A.38), which in more detail reads $$i\mathrm{\Delta }_{nm}^{RL}(t,t_2)=i\mathrm{\Delta }_{nm}^{LL}(t,t_2)=\overline{T}\{\xi _n(t)\xi _m(t_2)\}=\xi _m(t_2)\xi _n(t),$$ (A.41) in which we have used the fact that $`t>t_2`$. This, together with the first of Eqs. (A.37), tells us that $$i\mathrm{\Delta }_{nm}^{RL}(t_1,t_2)=\xi _m(t_2)\xi _n(t_1),$$ (A.42) which is the same propagator for internal lines connecting right vertices with left vertices that we found in the operator formalism. 5. Tree Graphs and Classical Solutions. We will now verify the remark made in Section I, that the usual approach to the calculation of non-Gaussian correlations, of solving the classical field equations beyond the linear approximation, simply corresponds to the calculation of tree diagrams in the “in-in” formalism. This is a well-known result<sup>13</sup> in the usual applications of quantum field theory, but some modifications in the usual argument are needed in the “in-in” formalism, in which the vacuum persistence functional is always unity whether or not we add a current term to the Lagrangian. We begin by introducing a generating functional $`W[j,t,g]`$ for correlation functions of fields at a fixed time $`t`$: $$e^{W[J,t,g]/g}\mathrm{vac},\mathrm{in}\left\|e^{\frac{1}{g}_a{\scriptscriptstyle d^3x\delta \varphi _a(𝐱,t)J_a(𝐱)}}\right\|\mathrm{vac},\mathrm{in}_g,$$ (A.43) where $`J_a`$ is an arbitrary current, and $`g`$ a real parameter, with the subscript $`g`$ indicating that the expectation value is to be calculated using a Lagrangian density multiplied with a factor $`1/g`$. (This is different from the usual definition of the effective action, because here we are not introducing the current into the Lagrangian.) The quantity of physical interest is of course $`W[J,t,1]`$, from which expectation values of all products of fields can be found by expanding in powers of the current. Using Eq. (A.27), we can calculate $`W`$ as the path integral $`e^{W[J,t,g]/g}={\displaystyle \delta \varphi _L\delta \pi _L\delta \varphi _R\delta \pi _R}`$ $`\times \mathrm{exp}\left(i{\displaystyle _{\mathrm{}}^t}𝑑t^{}{\displaystyle \frac{1}{g}}\stackrel{~}{L}[\delta \varphi _L,\delta \pi _L;t^{}]\right)`$ $`\times \mathrm{exp}\left(+i{\displaystyle _{\mathrm{}}^t}𝑑t^{}{\displaystyle \frac{1}{g}}\stackrel{~}{L}[\delta \varphi _R,\delta \pi _R;t^{}]\right)`$ $`\times {\displaystyle }\delta [\varphi _L(t)\delta \varphi _R(t)]{\displaystyle }\delta [\delta \pi _L(t)\delta \pi _R(t)]`$ $`\times e^{\frac{1}{g}_a{\scriptscriptstyle d^3x\delta \varphi _a(𝐱,t)J_a(𝐱)}}\mathrm{}`$ $`\times \mathrm{\Psi }_{\mathrm{vac}}[\delta \varphi _L(\mathrm{})]\mathrm{\Psi }_{\mathrm{vac}}[\delta \varphi _R(\mathrm{})]`$ (A.44) The usual power-counting arguments<sup>13</sup> show that the $`L`$ loop contribution to $`W[J,t,g]`$ has a $`g`$-dependence given by a factor $`g^L`$. For $`g0`$, $`W`$ is thus given by the sum of all tree graphs. The integrals over $`\delta \varphi _L`$, $`\delta \pi _L`$, $`\delta \varphi _L`$, $`\delta \pi _L`$ are dominated in the limit $`g0`$ by fields where $`\stackrel{~}{L}`$ is stationary, i.e., where $$\delta \varphi _L=\delta \varphi _R=\delta \varphi ^{\mathrm{classical}}$$ $$\delta \pi _L=\delta \pi _R=\delta \pi ^{\mathrm{classical}}$$ with $`\delta \varphi ^{\mathrm{classical}}`$ and $`\delta \pi ^{\mathrm{classical}}`$ the solutions of the classical field equations with the initial conditions that the fields go to free fields such as (14)–(16) satisfying the initial conditions (20) at $`t\mathrm{}`$. Since the $`L`$ and $`R`$ fields take the same values at this stationary point, the action integrals cancel, and we conclude that $$\left[W[J,t,1]\right]_{\mathrm{zero}\mathrm{loops}}=\underset{a}{}d^3x\delta \varphi _a^{\mathrm{classical}}(𝐱,t)J_a(𝐱).$$ (A.45) Expanding in powers of the current, this shows that in the tree approximation the expectation value of any product of fields is to be calculated by taking the product of the fields obtained by solving the non-linear classical field equations with suitable free-field initial conditions, as was to be proved. REFERENCES 1. For a review, see N. Bartolo, E. Komatsu, S. Matarrese, and A. Riotto, astro-ph/0406398. 2. J. Maldacena, JHEP 0305, 013 (2003) (astro-ph/0210603). For other work on this problem, see A. Gangui, F. Lucchin, S. Matarrese,and S. Mollerach, Astrophys. J. 430, 447 (1994) (astro-ph/9312033); P. Creminelli, astro-ph/0306122; P. Creminelli and M. Zaldarriaga, astro-ph/0407059; G. I. Rigopoulos, E.P.S. Shellard, and B.J.W. van Tent, astro-ph/0410486; and ref. 3. 3. J. Schwinger, Proc. Nat. Acad. Sci. US 46, 1401 (1961). Also see L. V. Keldysh, Soviet Physics JETP 20, 1018 (1965); B. DeWitt, The Global Approach to Quantum Field Theory (Clarendon Press, Oxford, 2003): Sec. 31. This formalism has been applied to cosmology by E. Calzetta and B. L. Hu, Phys. Rev. D 35, 495 (1987); M. Morikawa, Prog. Theor. Phys. 93, 685 (1995); N. C. Tsamis and R. Woodard, Ann. Phys. 238, 1 (1995); 253, 1 (1997); N. C. Tsamis and R. Woodard, Phys. Lett. B426, 21 (1998); V. K. Onemli and R. P. Woodard, Class. Quant. Grav. 19, 407 (2002); T. Prokopec, O. Tornkvist, and R. P. Woodard, Ann. Phys. 303, 251 (2003); T. Prokopec and R. P. Woodard, JHEP 0310, 059 (2003); T. Brunier, V.K. Onemli, and R. P. Woodard, Class. Quant. Grav. 22, 59 (2005), but not (as far as I know) to the problem of calculating cosmological correlation functions. 4. F. Bernardeau, T. Brunier, and J-P. Uzam, Phys. Rev. D 69, 063520 (2004). 5. The constancy of this quantity outside the horizon has been used in various special cases by J. M. Bardeen, Phys. Rev. D22, 1882 (1980); D. H. Lyth, Phys. Rev. D31, 1792 (1985). For reviews, see J. Bardeen, in Cosmology and Particle Physics, eds. Li-zhi Fang and A. Zee (Gordon & Breach, New York, 1988); A. R. Liddle and D. H. Lyth, Cosmological Inflation and Large Scale Structure (Cambridge University Press, Cambridge, UK, 2000). 6. R. S. Arnowitt, S. Deser, and C. W. Misner, in Gravitation: An Introduction to Current Research, ed. L. Witten (Wiley, New York, 1962): 227. This classic article is now available as gr-qc/0405109. 7. V. S. Mukhanov, H. A. Feldman, and R. H. Brandenberger, Physics Reports 215, 203 (1992); E. D. Stewart and D. H. Lyth, Phys. Lett. B 302, 171 (1993). 8. N. C. Tsamis and R. Woodard, ref. 3. 9. A. Vilenkin and L. H. Ford, Phys. Rev. D. 26, 1231 (1982); A. Vilenkin, Nucl. Phys. B226, 527 (1983). 10. R. Kubo, J. Math. Phys. 4, 174 (1963). 11. See, e.g., S. Weinberg, The Quantum Theory of Fields – Volume I (Cambridge, 1995): Sec. 9.1. 12. ibid., Sec. 9.2. 13. S. Coleman, in Aspects of Symmetry (Cambridge University Press, Cambridge, 1985): pp 139–142.
warning/0506/astro-ph0506622.html	ar5iv	text	# The Ital-FLAMES survey of the Sagittarius dwarf Spheroidal galaxy. I. Chemical abundances of bright RGB stars Based on observations obtained with FLAMES at VLT Kueyen 8.2m telescope in the program 71.B-0146. ## 1 Introduction The Local Group (LG) is a heterogeneous environment. Galaxies in the LG show a variety of characteristics (e.g. mass, morphology, gas content) and are evolving under different conditions (e.g. in isolation, on strong dynamical interaction). Thererore, in principle, they could teach us about galaxy evolution as much as globular clusters did concerning stellar evolution. Chemical abundances and abundance ratios are key ingredients to study the star formation histories of stellar systems. The modern generation of spectrographs mounted on 8-10 m class telescopes allows to investigate the chemical composition and dynamics of bright stars in LG galaxies but only a handful of stars have been studied so far (Tolstoy et al., 2004, 2003; Shetrone et al., 2003; Bonifacio et al., 2000, 2004; Fulbright et al., 2004; Geisler et al., 2005; Shetrone et al., 1998, 2001). The commonly accepted paradigm (White & Rees, 1978) predicts the formation of large galaxies from the hierarchical assembly of small fragments similar to the LG dwarf spheroidals (dSphs). In this framework, the comparison between the chemical composition of the Milky Way (MW) and LG dSph stars is a first local testbed for the hierarchical merging model. The chemical composition of LG stars turned out to be remarkably different from that of MW stars of comparable metallicities. In particular, LG stars show $`\alpha `$ element abundance ratios systematically under-abundant with respect to MW stars (see, for instance, Venn et al., 2004; Bonifacio et al., 2004). The interpretation of this empirical evidence is controversial. Is the hierarchical merging a minor process in the assembly of the MW? Or were the fragments from which the MW formed at early times different from the nowadays recognizable dSphs? The chemical difference between MW and LG stars may reflect an environmental difference between dwarfs accreted at early times (galaxies near the bottom of the pre-MW potential well – dense environment) and the surviving dwarfs (galaxies far from the bottom of the pre-MW potential well – loose environment, but see Robertson et al., 2005; Bullock & Johnston, 2004). The Sagittarius dSph (hereafter Sgr, Ibata, Gilmore, & Irwin, 1994) is a LG galaxy currently experiencing strong and disruptive tidal interactions with the MW (Ibata, Gilmore, & Irwin, 1995; Ibata et al., 1997; Majewski et al., 2003). Therefore, it may provide clues on the influence of dynamical interactions on the chemical evolution of dwarf galaxies. It is well-known that the complex stellar content of Sgr (see Monaco et al., 2002, 2003, 2005, and references therein) is largely dominated by a population of old-intermediate age stars ($``$6 Gyr, see, e.g. Bellazzini et al., 1999; Layden & Sarajedini, 2000; Monaco et al., 2002). However, some concerns have been raised on mean metallicity estimates obtained for this population from spectroscopic and photometric works (see, e.g. Mateo et al., 1995; Bonifacio et al., 2000; Cole, 2001; Monaco et al., 2002; Bonifacio et al., 2004). The paper is devoted to the assessment of the mean chemical properties of the Sgr dominant population. We present Fe, Mg, Ca and Ti abundances for a selected sample of stars belonging to this population. In a companion paper (Bonifacio et al., in preparation) we deal with the issue of the Sgr metallicity distribution. \[Fe/H\] and \[$`\alpha `$/Fe\] abundances as well as the trends in the \[Fe/H\] vs \[$`\alpha `$/Fe\] plane constrain the chemical evolution which led to the formation of the Sgr dominant population (Lanfranchi & Matteucci, 2003). Moreover, mean \[Fe/H\] and \[$`\alpha `$/Fe\] values are key ingredients to derive reliable age estimates from the color-magnitude diagrams. The paper is organized as follows. In §2 we describe the target selection and the obtained data. In §3 we describe the procedures followed to fix the atmospheric parameters and the chemical analysis. In §4 we compare the results obtained with previous works and in §5 we discuss our findings. ## 2 Observations ### 2.1 Target selection, Data and Equivalent Widths As part of the guaranteed time awarded to the Ital-FLAMES consortium, more than 400 stars were observed in the Sgr dSph (Bonifacio et al., 2005; Zaggia et al., 2004) from May the 23th to 27th, 2003, using the FLAMES facility mounted on the VLT (Pasquini et al., 2000). Details on the observations are given in Zaggia et al. (2004). FLAMES allows to observe 132 targets in one shot using the intermediate-low resolution spectrograph GIRAFFE, plus 8 additional targets using the red arm of the high resolution specrograph UVES. In this paper we present the results obtained from the UVES spectra. It is important to recall that a large number of Milky Way foreground stars are present along the Sgr line of sight. In order to optimize the Sgr star detection rate, the target selection for the UVES fibres was performed using the infrared 2 MASS<sup>1</sup><sup>1</sup>1See http://www.ipac.caltech.edu/2mass color magnitude diagram (CMD). In fact, in the infrared plane, the upper Sgr red giant branch (RGB) stands out very clearly from the contaminating MW field (see, e.g., Cole, 2001). This also allows a thorough sampling of the Sgr dominant population. In Fig. 1 we plotted the 2 MASS (K; J-K<sub>S</sub>) CMD for a 1 square degree area centred on the globular cluster M 54. The heavy continuous line is the selection box. Target stars are plotted as large filled circles. A similar target selection already proven to be very effective in detecting stars belonging to the Sgr Stream (Majewski et al., 2004). Target stars are marked as large symbols in the optical CMD plotted in Fig. 2 (Monaco et al., 2002). In Table 1 we report equatorial (J 2000.0) coordinates and V and I magnitudes for the target stars. The coordinates in the J2000.0 absolute astrometric system for both UVES and GIRAFFE samples were obtained with a procedure already described in other papers (see, for example, Ferraro et al., 2001). The new astrometric Guide Star Catalogue (GSC II) recently released and now available on the web<sup>2</sup><sup>2</sup>2See http://www-gsss.stsci.edu/gsc/gsc2/GSC2home.htm was used as reference. In order to derive an astrometric solution we used a program specifically developed at Bologna Observatory (Montegriffo et al., in preparation). As a result of the entire procedure, rms residuals of $``$0.15 arcsec, both in RA and Dec, were obtained. The quality of the astrometry was confirmed by the successful centreing of the fibres. We performed the analysis on the spectra reduced with the UVES ESO-MIDAS<sup>3</sup><sup>3</sup>3ESO-MIDAS is the acronym for the European Southern Observatory Munich Image Data Analysis System which is developed and maintained by the European Southern Observatory. http://www.eso.org/projects/esomidas/ pipeline. For each pointing, 7 fibres were centred on the target stars while one fibre was used to measure the sky spectrum. Different spectra of the same star were coadded and the resulting signal to noise ratio (S/N) ranges from 14 to 43 at 653 nm (see Table 2). UVES spectra have a resolution of R$``$43000 and cover the range between 480 nm and 680 nm. Equivalent widths (EW) were measured on the spectra using the standard IRAF<sup>4</sup><sup>4</sup>4IRAF is distributed by the National Optical Astronomy Observatories, which is operated by the association of Universities for Research in Astronomy, Inc., under contract with the National Science Foundation. task splot. The Fe, Ca, Mg and Ti line lists as well as the adopted atomic parameters and the measured EW are reported in Table 5. A different iron line list (see Table 6) was adopted for star #3800319 due to the relatively high temperature and gravity of this star in comparison with the other stars in the sample. We analysed interactively the spectral lines. For each line the fit has been visually inspected and adjusted until reaching a satisfying solution. ### 2.2 Radial velocities and the contaminating Milky Way field Radial velocities (see Table 1) were obtained by cross-correlating the observed spectra with a rest frame laboratory line list using the recently released software DAOSPEC<sup>5</sup><sup>5</sup>5See http://cadcwww.hia.nrc.ca/stetson/daospec (Stetson and Pancino, in preparation). The final radial velocities and relative errors were computed using about 150 lines for each star. Geocentric observed radial velocities were corrected to heliocentric velocities using the IRAF task rvcorrect. The DAOSPEC code has the capability to measure the line EWs. In our case, however, we used DAOSPEC only to measure the radial velocities of the target stars while we used the IRAF task splot to measure EWs for homogeneity with our previous works on Sgr stars (Bonifacio et al., 2000, 2004). As a check, the radial velocities of a few stars have also been measured using the fxcor IRAF task for Fourier cross correlation. The radial velocities obtained using DAOSPEC and fxcor are identical, within the errors. All but one (#3600127, v<sub>helio</sub>=-127.6 km/s, open circle in Fig. 2) of the 24 observed stars are indeed Sgr radial velocity members lying within $``$2$`\sigma `$ of the systemic velocity as measured by Ibata et al. (1997). In Fig. 3 we plotted the velocity distribution of the 23 Sgr radial velocity members. The mean velocity ($`<`$v$`{}_{r}{}^{}>`$=143.08$`\pm `$3.2 km s$`{}_{}{}^{1})`$<sup>6</sup><sup>6</sup>6The quoted 3.2 km s<sup>-1</sup> error has been estimated employing a bootstrap technique. and the velocity dispersion of the sample ($`\sigma `$=11.17 km s<sup>-1</sup>) are in good agreement with the values derived by Ibata, Gilmore, & Irwin (1995) and Ibata et al. (1997). The MW model of Robin et al. (2003, hereafter R03) predicts that in the M 54 line of sight 2% of stars have v$`{}_{r}{}^{}>`$100 km s<sup>-1</sup>, if we consider only stars lying in the same (V, V-K)<sup>7</sup><sup>7</sup>7The R03 model does not provide the (J-K) color, therefore we define as selection box in the (V, V-K) plane the region which encloses all the target stars. selection box of the UVES sample. However, the model predicts only a $``$4% of giant stars (log g$`<`$4) in the selection box and none of them with v$`{}_{r}{}^{}>`$100 km s<sup>-1</sup>. We checked carefully the 24 stars in the sample and we are confident that all of them are indeed red giant. Therefore, even if the R03 model provides only an approximate description of the MW, there is no reason to expect any MW star among the 23 Sgr radial velocity members in the sample. ### 2.3 M-giants showing TiO molecular bands in the spectra The coolest (i.e. the reddest) four stars (#2300168, #3600181, #3700055, #4207391, plus symbols at V-I$`>`$2.0 in Fig. 2) have effective temperatures around 3600 K and very strong titanium oxide bands (TiO, see Selvelli & Bonifacio, 2000; Valenti et al., 1998) in the spectra (see Fig. 4). The presence of the TiO bands confirm these stars as M-giants. Such strong molecular bands prevent from a safe derivation of the equivalent widths. Therefore, we do not present the chemical analysis for these stars. In addition, stars #3600073, #3700178, #3800366, #4207953 (plus symbols at V-I$`<`$2.0 in Fig. 2) show weak but clearly recognizable TiO bands. For these stars we provide only a tentative analysis and the derived abundances will not be discussed. We plan to provide a detailed chemical analysis for these 8 stars by performing spectral synthesis including also the TiO molecular bands. In the CMD in Fig. 2 star #2300127 lies exactly in the region occupied by stars with TiO bands in their spectra. Yet this star does not present any band. The lack of the TiO molecular bands may be due to the relatively weak Fe and Ti content of this star (\[Fe/H\]=-0.81, \[Ti/Fe\]=-0.17, see Table 3). ## 3 Atmospheric Parameters and Chemical analysis The UVES spectra of the 19 stars for which the chemical analysis was performed (including also stars having weak TiO bands) are plotted in Fig. 5. ### 3.1 Effective temperatures and surface gravities The effective temperatures for the target stars (see Table 1) were derived from the (V-I) color assuming a reddening of E(V-I)=0.18 (Layden & Sarajedini, 2000) and using the calibration of Alonso, Arribas, & Martínez-Roger (1999). We used the Girardi et al. (2002) theoretical isochrones, along with E(V-I)=0.18 and (m-M)<sub>0</sub>=17.10 (Layden & Sarajedini, 2000; Monaco et al., 2004) as reddening<sup>8</sup><sup>8</sup>8We assumed the same reddening for all the stars in the sample. Inspection of the Schlegel et al. (1998) reddening maps provide strong indications that there is no serious variability of extinction in the considered field (standard deviation of the reddening value: $`\sigma _{E(BV)}`$=0.03, Monaco et al., 2004). and distance modulus, in order to estimate the gravity of the program stars. In particular, we used a (Z=0.001; Age =14.13 Gyr) isochrone for stars #3800199 #3800204 #3800319 and a (Z=0.008; Age =6.31 Gyr) isochrone for all the other stars (continuous lines in Fig. 2). These two isochrones fit into the range covered by the target stars on the CMD and the age and metallicity used are also compatible to what expected from previous works (see Monaco et al., 2002; Layden & Sarajedini, 2000; Brown, Wallerstein, & Gonzalez, 1999; Bonifacio et al., 2004). ### 3.2 Model atmosphere and Microturbulent velocities For each star we computed a plane parallel model atmosphere using version 9 of the ATLAS code (Kurucz, 1993) with the above atmospheric parameters. Abundances were derived from EWs using the WIDTH code (Kurucz, 1993). Microturbulent velocities ($`\xi `$) were determined minimizing the dependence of the iron abundance from the EW, among the set of iron lines measured for each star. In Fig. 6 we plotted $`\xi `$ as a function of the adopted gravity for the stars studied by Ivans et al. (2001, bottom panel, hereafter I01), Shetrone et al. (2003, middle panel, hereafter S03) and for stars in our sample (top panel). A clear trend is present in both the I01 and S03 samples. The same trend is present also in our sample, albeit with a larger scatter. Continuous lines are least square fits to the data points. In the case of our sample the fit was obtained excluding the points having the highest and lowest $`\xi `$ (2.7 and 0.6 km s<sup>-1</sup>) and the point with the lowest surface gravity (log g=0.41). As can be seen, the three fitting lines are very similar to each other. A weak dependence of the $`\xi `$ from the effective temperature was found and it can be safely neglected as a first order approximation. For stars #2300127, #2300215 and #3800319 (filled circles in fig. 6) the $`\xi `$ is 2.7, 2.5 and 0.6 km s<sup>-1</sup>, respectively, i.e. more than 2-$`\sigma `$ far from the fitting relation. When working with low S/N, highly crowded spectra, it is difficult to measure weak Fe lines accurately. This may lead to incorrect $`\xi `$. Thus, for stars #2300127, #2300215 and #3800319 we adopted the value obtained from the fitting relation $`\xi `$=-0.35$`\times `$log g+2.29, i.e. $`\xi `$=2.0, 1.9 and 1.9 km s<sup>-1</sup>, respectively. ### 3.3 Chemical Abundances The atmospheric parameters (T<sub>eff</sub>, log g, $`\xi `$ and the assumed global metallicity \[M/H\]) adopted for the program stars are reported in Table 1. The chemical abundances obtained for each line are reported in Tables 5 and 6. The mean and standard deviation of such abundances are reported in Tables 2 and 3 (as \[X/H\] abundances in the latter case) for each chemical species. In Table 2 we also reported the number of lines used to obtain the mean abundance for each species. The line scatter reported in Table 2 should be representative of the statistical error arising from the noise in the spectra and from uncertainties in the measurement of the equivalent widths. Under the assumption that each line provides an independent measure of the abundance, the error in the mean abundances should be obtained by dividing the line scatter by $`\sqrt{n}`$ (where $`n`$ is the number of measured lines) and by adding to this figure the errors arising from the uncertainties in the atmospheric parameters. In Table 4 we report these latter errors in the case of star #3800318, taken as representative of the whole sample. In Fig. 7 we plotted the metallicity distribution obtained. Our sample spans a rather large metallicity range (-1.52$``$\[Fe/H\]$``$-0.17). The distribution peaks around \[Fe/H\]$``$-0.4 and presents an extended metal poor tail<sup>9</sup><sup>9</sup>9Preliminary results obtained from the GIRAFFE sample show that such tail extends at least down to \[Fe/H\]$`<`$-2.5 (Zaggia et al., 2004; Bonifacio et al., 2005). In particular, considering only stars more metal rich than \[Fe/H\]$``$-1, which should be representative of the Sgr dominant population, we obtain a mean value of $`<`$\[Fe/H\]$`>`$=-0.41$`\pm `$0.20. In Fig. 8 we plotted the \[Ti/Fe\], \[Ca/Fe\] and \[Mg/Fe\] ratios (from top to bottom panel) for the program stars as a function of the \[Fe/H\] abundance. The 5 M 54 stars studied by Brown, Wallerstein, & Gonzalez (1999, hereafter B99) are plotted as large open stars. Assuming $`<[\alpha /Fe]>=\frac{[Mg/Fe]+[Ca/Fe]}{2}`$, we also obtain a mean value of $`<`$\[$`\alpha `$/Fe\]$`>`$=-0.17$`\pm `$0.07 for the dominant population. Following Salaris, Chieffi & Straniero (1993), these values correspond to a global metallicity<sup>10</sup><sup>10</sup>10The “global metallicity” is defined as: \[M/H\]=\[Fe/H\]+log(0.638$`\times `$10<sup>\[α/Fe\]</sup>+0.362) of \[M/H\]=-0.51, which is in good agreement with the recent photometric estimate by Monaco et al. (2002). ### 3.4 Notes on Metal poor stars: #3800199, #3800204,#3800319 The three most metal poor stars (#3800199 #3800204 and #3800319) occupy in the optical CMD (see Fig. 2) positions compatible with the M 54 RGB (which is roughly represented by the bluer isochrone in the plot). The most metal poor star (#3800204, \[Fe/H\]=-1.52) lies very near to the M 54 center ($``$1$`\mathrm{}`$) and its chemical abundances (Fig. 8) are identical to those of the M 54 stars studied by B99. Therefore, it seems quite likely that this star does indeed belong to M 54. Star #3800319 (\[Fe/H\]=-1.37) is also quite near ($``$1$`\mathrm{}`$.4) to the cluster center but its chemical composition is only marginally compatible with M 54 and it will be considered a Sgr field star. However, we note that Layden & Sarajedini (2000) claimed a metallicity dispersion of $``$0.16 dex for M 54 from the width of the red giant branch. Star #3800199 (\[Fe/H\]=-1.10) is placed at 3$`\mathrm{}`$.2 from the cluster center (which corresponds to $``$7 half light radii, Trager, King, & Djorgovski, 1995) and is significantly more metal rich than M 54 (\[Fe/H\]$``$-1.55, B99). Therefore, we consider star #3800199 part of the Sgr galaxy field. ## 4 Comparison with previous works Beside the present work, chemical abundances have been presented for Sgr RGB stars by Bonifacio et al. (2000, 2004, 2 and 10 stars, respectively) and by Smecker–Hane & McWilliam (2002, hereafter S02, 14 stars). Bonifacio et al. (2004, hereafter B04) also considered the two stars studied in Bonifacio et al. (2000) obtaining a final sample of 12 stars. In Fig. 9 we plotted the \[$`\alpha `$/Fe\] as a function of the iron abundance for the stars in the 3 samples. Stars in our sample are plotted as filled circles, while stars in the B04 and S02 sample are plotted as empty squares and empty triangles, respectively. The 5 M 54 stars studied by B99 are marked as large open stars. The $`\alpha `$ element abundance ratio is defined as $`[\alpha /Fe]=\frac{[Mg/Fe]+[Ca/Fe]}{2}`$ for stars in our sample and in the B04 and B99 samples, while it is defined as $`[\alpha /Fe]=\frac{[Si/Fe]+[Ca/Fe]+[Ti/Fe]}{3}`$ for stars in the S02 sample<sup>11</sup><sup>11</sup>11S02 do not provide abundances for each species but only mean values.. Stars in the S02 sample range from \[Fe/H\]$``$-1.6 to \[Fe/H\]$``$0. In particular, 3 stars in their sample have \[Fe/H\]$`<`$-1 and 11 stars are in the range -0.7$`÷`$0.0. This latter sub-sample has a mean metallicity and $`\alpha `$ element abundance ratio of: $`<`$\[Fe/H\]$`>`$=-0.36$`\pm `$0.19 and $`<`$\[$`\alpha `$/Fe\]$`>`$=+0.01$`\pm `$0.04. However, it is important to note that these values should not be considered as representative of the dominant population, since their target selection has been biased toward stars with metallicities within 0.5 dex of the solar value based on previously obtained approximate metallicities (McWilliam & Smecker-Hane, 2004). The metallicity range of stars in the B04 sample, on the other hand, is -0.83$``$\[Fe/H\]$`<`$+0.09. Therefore, it extends toward slightly higher metallicity with respect to the S02 sample, but it lacks of metal poor stars. The mean metallicity and $`\alpha `$ element abundance ratio of the B04 sample are: $`<`$\[Fe/H\]$`>`$=-0.23$`\pm `$0.26 and $`<`$\[$`\alpha `$/Fe\]$`>`$=-0.20$`\pm `$0.06. The mean iron abundance obtained in this paper (\[Fe/H\]=-0.41) is similar to that of the S02 and B04 samples. The 0.18 dex difference between the B04 mean iron abundance and our figure would be also a little bit lowered (by 0.06$`÷`$0.09 dex) by taking into account the different assumption about the reddening (B04 adopted E(V-I)=0.22 from Marconi et al., 1998). The different target selection criterion adopted by B04 may also be responsible for the residual difference in the mean iron abundance ($``$0.1 dex), which is, nevertheless, well inside the involved errors. The $`<`$\[$`\alpha `$/Fe\]$`>`$ ratio obtained by B04 is very similar to our value ($`<[\alpha /Fe]>=\frac{[Mg/Fe]+[Ca/Fe]}{2}`$=-0.17). S02 evaluate the $`\alpha `$ element abundance ratio as $`[\alpha /Fe]=\frac{[Si/Fe]+[Ca/Fe]+[Ti/Fe]}{3}`$. Considering $`[\alpha /Fe]=\frac{[Ca/Fe]+[Ti/Fe]}{2}`$, we obtain a $`<[\alpha /Fe]>`$ fairly similar to the S02 figure. The small residual difference ($``$0.1 dex higher in the S02 sample) may be partly ascribed to the \[Si/Fe\] abundances and, possibly, to a different set of lines and atomic parameters adopted in the chemical analysis. Unfortunately, S02 neither provide abundances for each species nor the atomic data and the adopted line list and this hypothesis cannot be checked further. Finally, as already stressed in section 3.4, we remark that the Fe, Mg, Ca and Ti abundances of star #3800204 are consistent with the results obtained by B99 for M 54 stars. ## 5 Discussion and Conclusions The main purpose of this paper was to study the chemical composition of the dominant population of the Sgr dSph galaxy. We selected 24 target stars using the 2 MASS infrared CMD, where the upper RGB of Sgr is well separated from the MW field. Target stars have been observed using the red arm of the high resolution spectrograph FLAMES-UVES. We reported radial velocities for these 24 stars and all but one are Sgr radial velocity members. Eight stars show strong or visible TiO bands. For stars with weak TiO bands we present a tentative chemical analysis while we do not present any chemical analysis for stars presenting strong TiO bands in the spectra. For the remaining 15 stars, we reported Fe, Mg, Ca and Ti chemical abundances. This is the largest sample of high resolution spectra analyzed so far for stars in the Sgr dSph galaxy, and the only sample thoroughly representative of the Sgr dominant population. The metallicity ranges from \[Fe/H\]=-1.52 to \[Fe/H\]=-0.17. Three stars have \[Fe/H\]$`<`$-1 and the most metal poor of them (#3800204) can be reasonably considered M 54 member. The mean iron content of stars with \[Fe/H\]$`>`$-1 (i.e. the Sgr dominant population) is $`<`$\[Fe/H\]$`>`$=-0.41$`\pm `$0.20, with a mean $`\alpha `$ element abundance ratio $`<`$\[$`\alpha `$/Fe\]$`>`$=-0.17$`\pm `$0.07. These figures lead to a global metallicity \[M/H\]=-0.51 which is in close agreement with the most recent photometric estimates obtained for the Sgr dominant population (Monaco et al., 2002). In order to obtain a more statistically significant sample, we now join the B04 and our samples. In Fig. 10 we plotted in the \[$`\alpha `$/Fe\] vs \[Fe/H\] plane the mean points obtained for Sgr from this larger sample of Sgr stars as filled circles. For $`0.65<`$\[Fe/H\]$`<0.1`$, filled circles represent running means with 0.20 dex as bin and 0.1 dex as step. For stars having $`1.0<`$\[Fe/H\]$`<0.65`$ and $`1.5<`$\[Fe/H\]$`<1`$ (i.e. excluding star #3800204 which has been tagged as M 54 member) filled circles are straight means of the \[Fe/H\] and \[$`\alpha `$/Fe\] with the corresponding standard deviations as errorbars. A weak, but clearly recognizable trend between the $`\alpha `$ element abundance ratio and the mean iron abundance exists at high metallicity. Such a trend waits to be confirmed from a much more extended sample such as that obtained using the FLAMES-GIRAFFE multifibre spectrograph which is currently under analysis. For \[Fe/H\]$`<`$-1, a sudden increase of the \[$`\alpha `$/Fe\] is apparent. The mean \[$`\alpha `$/Fe\] at low metallicities is consistent with the values observed in MW stars (crosses in Fig. 10, from Venn et al., 2004) of comparable metallicities and somewhat higher with respect to stars in the LG galaxies (asterisks in the figure, from Venn et al., 2004). Therefore, metal poor stars lost in early passages which now are not recognizable as Sgr tidal debris (Helmi, 2004), would be part of the typical content of the MW Halo and impossible to tag as an accreted component from the chemical composition. The three metal poor stars in the S02 sample are compatible with MW stars as well. This occurrence led the authors to suggest that the upper mass end of the Sgr initial mass function (IMF) should not be significantly different from the MW one. The level of \[$`\alpha `$/Fe\] which characterizes a galaxy at low metallicities may indeed give information on the IMF of the galaxy at that time (see McWilliam, 1997, and references therein), since the amount of $`\alpha `$ elements and iron produced by a Type II SN is a function of the mass of the SN progenitor. Although this is true in principle, in practice this information may not be presently extracted. In fact the ratio of the $`\alpha `$ elements and iron produced by a Type II SN is also a sensitive function of the “mass cut”, i.e. the mass coordinate which separates the material of the SN which “falls back” on the SN remnant from the material which is ejected. The deeper the mass cut, the more iron-peak elements are ejected, thus lowering the overall \[$`\alpha `$/Fe\]. Current SN models are unable to determine the mass cut in a self consistent way or from first principles: the mass cut is always assumed. We do not have either any indication on whether the mass cut is in any sense “Universal” or if it may vary e.g. depending on the mass of the star or on its metallicity. With this state of affairs, any inference on the IMF from the level of the \[$`\alpha `$/Fe\] ratio of a galaxy would be highly uncertain. The metal-rich Sgr stars lie on the extension to high metallicity of the pattern followed by stars in LG galaxies and below MW stars. A low \[$`\alpha `$/Fe\] at high metallicity is traditionally interpreted as evidence for a slow or bursting star formation rate (S02, B04, Marconi et al., 1994). On the contrary, in order to reproduce the Sgr \[$`\alpha `$/Fe\] ratios, Lanfranchi & Matteucci (2003) required a high star formation rate. However, they constrained their model using preliminary abundances presented by Smecker-Hane & Mc William (1999). The somewhat lower \[$`\alpha `$/Fe\] values obtained here and in S02 and B04 should be in better agreement with a lower star formation rate. Our data suggest that Sgr had a different chemical evolution from both the MW and the LG galaxies (see Fig. 10). A different chemical evolution for Sgr with respect to the other LG galaxies is expected, since Sgr experienced strong and disruptive dynamical interactions with the MW. Such interactions are witnessed by the Sgr tidal tails studied by Majewski et al. (2003) and are expected to trigger star formation activity (see, for instance, Kravtsov et al., 2004; Mayer et al., 2001; Zaritsky & Harris, 2004). Finally, we note that the Large Magellanic Cloud (LMC) metallicity distribution strongly resambles the Sgr one. In fact, Cole et al. (2005) approximated the metallicity distribution of the LMC bar by two Gaussians having \[Fe/H\]=-0.37$`\pm `$0.15 and \[Fe/H\]=-1.08$`\pm `$0.46 and containing 89% and 11% of the stars, respectively. The same results hold also for the LMC disk (see Cole et al., 2005). Clearly, Sgr has the same mean metallicity of the LMC dominant population as well as the same fraction of metal poor stars (see Monaco et al., 2003). Such occurrence may suggest a similarity of the Sgr progenitor with the LMC. ###### Acknowledgements. Part of the data analysis has been performed using software developed by P. Montegriffo at the INAF - Osservatorio Astronomico di Bologna. This research was done with support from the Italian MIUR COFIN/PRIN grants 2002028935 and 2004025729. We are grateful to L. Girardi for useful comments and to G. Schiulaz for a careful reading of the manuscript. ## Appendix A Individual line data The following tables report the line list and adopted atomic parameters for the program stars. The measured equivalent width and the corresponding abundance obtained for each line are also reported.
warning/0506/nucl-th0506020.html	ar5iv	text	# 𝜂-meson in nuclear matter ## I Introduction The studies of meson-baryon interactions and the meson properties in nuclear medium are interesting subjects in nuclear physics. The pion-nucleon/pion-nucleus and kaon-nucleon/kaon-nucleus interactions have been much studied, both theoretically and experimentally. Due to the lack of eta beams, the $`\eta `$-nucleon/$`\eta `$-nucleus interaction is still not as clear as that of the pion-nucleon/pion-nucleus and kaon-nucleon/kaon-nucleus. Since the $`\eta `$-nucleus quasi-bound states were first predicted by Haider and Liu b1 and Li et al. b2 , when it was realized that the $`\eta `$-nucleon interaction is attractive, the study of the $`\eta `$-nucleus bound states has been one of the focuses in nuclear physics b4 ; b5 ; b6 ; b7 ; b8 ; b9 ; b10 ; b11 ; b12 . The key point for the study of $`\eta `$-nucleus bound states is the $`\eta `$ nuclear optical potential. There have been some works in this field. Waas and Weise studied the s-wave interactions of $`\eta `$-meson in nuclear medium, and got a potential $`U_\eta 20`$ MeV b3 . Chiang et alop gave $`U_\eta 34`$ MeV by assuming that the mass of the $`N^{}(1535)`$ did not change in the medium. Tsushima et al. predicted that the $`\eta `$-meson potential was typically $`60`$ MeV using QMC model efm . Inoue and Oset also obtained $`U_\eta 54`$ MeV with their model op1 . Obviously, there are model dependencies in describing the in-medium properties of $`\eta `$-meson. Therefore, further studies are needed. In this paper, firstly we deduce the $`\eta `$N interactions from chiral perturbation theory; then combining the relativistic mean-field theory for nucleon system, we will study the properties of eta meson in uniform nuclear matter. The relativistic mean field theory (RMF) is one of the most popular methods in modern nuclear physics. It has been successful in describing the properties of ordinary nuclei/nuclear matter and hyper-nuclei/nuclear matterrmf ; Lalazissis97prc55 . On the other hand, the chiral perturbation theory (ChPT) was first applied by Kaplan and Nelson to investigate the in-medium properties of (anti)kaons chiral01 . Some years later, an effective chiral Lagrangian in heavy-fermion formalism chiral1 was also introduced to study the kaon-nuclear/nucleon interactions or kaon condensation chiral ; CHL ; chiralpp . The advantage of using the heavy-fermion Lagrangian for chiral perturbation theory was clearly pointed out in Ref. chiral1 . Compared with the previous chiral perturbation theory chiral01 , the outstanding point in Refs. chiral ; CHL ; chiralpp is that additional next-to-leading-order terms, i.e., off-shell terms, are added to the Lagrangian. The additional terms are essential for a correct description of the KN interactions. The chiral perturbation theory also had been used in the study of $`\eta `$-meson in-medium properties in Ref. b3 ; op1 , where only the leading-order terms were kept in the calculations. Given that the higher order terms, e.g., off-shell terms, are important to the $`\eta `$N interactions, and they have not been included in the previous studies for the $`\eta `$N interactions with chiral perturbation theory, we have, in the present work, studied the $`\eta `$N interactions with the heavy-baryon chiral perturbation theory up to the next-to-leading-order terms. Combining the RMF for nuclear matter, we obtain the in-medium properties of $`\eta `$-meson. Comparing our results with the previous results (with only leading-order terms), we find that the next-to-leading-order terms are important to the calculations indeed. The $`\eta `$-nucleon sigma term is found to be 280 $`\pm `$ 130 MeV. The ratio of $`\eta `$-meson effective mass to its vacuum value is $`0.84\pm 0.015`$, while depth of the optical potential is $`(83\pm 5)`$ MeV, at the normal nuclear density. The large uncertainty in the sigma term $`\mathrm{\Sigma }_{\eta N}`$ does not affect the results significantly in low density region, varying by about 8 MeV at normal nuclear density. The paper is organized as follows. In the subsequent section, the effective chiral Lagrangian density we used is given, the effective Lagrangian for $`\eta `$N interactions is derived, and the coefficients for the sigma and off-shell terms are determined. Then, combining the RMF for nucleons, we obtain the $`\eta `$-meson energy, effective mass, and optical potential in nuclear matter in Sec. III. We present our results and discussion of the $`\eta `$-meson in-medium properties in Sec. IV. Finally a summary is given in Sec. V. ## II The $`\eta `$N interactions in Chiral perturbation theory ### II.1 The theory framework The interactions between pseudoscalar mesons (pion, kaon, and eta meson) and baryons (nucleons and hyperons) are described by the SU(3)$`{}_{\mathrm{L}}{}^{}\times `$SU(3)<sub>R</sub> chiral Lagrangian which can be written as $`_{chiral}=_\varphi +_{\varphi B}.`$ (1) $`_\varphi `$ is the mesonic term up to second chiral orderchiral01 , $`_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{4}}f^2\text{Tr}^\mu \mathrm{\Sigma }_\mu \mathrm{\Sigma }^{}`$ (2) $`+{\displaystyle \frac{1}{2}}f^2B_0\{\text{Tr}M_q(\mathrm{\Sigma }1)+\mathrm{h}.\mathrm{c}.\}.`$ The second piece of the Lagrangian in Eq.(1), $`_{\varphi B}`$, describes the meson-baryon interactions, and reads at lowest orderchiral01 $`_{\varphi B}^{(1)}`$ $`=`$ $`\text{Tr}\overline{B}(i\gamma ^\mu _\mu m_\mathrm{B})B+i\text{Tr}\overline{B}\gamma ^\mu [V_\mu ,B]`$ (3) $`+D\text{Tr}\overline{B}\gamma ^\mu \gamma ^5\{A_\mu ,B\}+F\text{Tr}\overline{B}\gamma ^\mu \gamma ^5[A_\mu ,B],`$ The next-to-leading order chiral Lagrangian for s-wave meson-baryon interactions reads chiral $`_{\varphi B}^{(2)}`$ $`=`$ $`a_1\text{Tr}\overline{B}(\xi M_q\xi +\text{h.c.})B+a_2\text{Tr}\overline{B}B(\xi M_q\xi +\text{h.c.})`$ (4) $`+a_3\text{Tr}\overline{B}B\text{Tr}(M_q\mathrm{\Sigma }+\text{h.c.})+d_1\text{Tr}\overline{B}A^2B`$ $`+d_2\text{Tr}\overline{B}(vA)^2B+d_3\text{Tr}\overline{B}BA^2`$ $`+d_4\text{Tr}\overline{B}B(vA)^2+d_5\text{Tr}\overline{B}B\text{Tr}A^2`$ $`+d_6\text{Tr}\overline{B}B\text{Tr}(vA)^2+d_7\text{Tr}\overline{B}A_\mu \text{Tr}A^\mu B`$ $`+d_8\text{Tr}\overline{B}(vA)\text{Tr}(vA)B+d_9\text{Tr}\overline{B}A_\mu BA^\mu `$ $`+d_{10}\text{Tr}\overline{B}(vA)B(vA),`$ In the above equations, $`M_q=\text{diag}\{m_q,m_q,m_s\}`$ is the current quark mass matrix, $`B_0`$ relates to the order parameter of spontaneously broken chiral symmetry, the constants $`D`$ and $`F`$ are the axial vector couplings whose values can be extracted from the empirical semileptonic hyperon decays, the pseudoscalar meson decay constants are equal in the SU(3)<sub>V</sub> limit, and denoted by $`f=f_\pi 93`$ MeV, $`v_\mu `$ is the four-velocity of the heavy baryon (with $`v^2`$=1), and $`\mathrm{\Sigma }=\xi ^2=\mathrm{exp}(i\sqrt{2}\mathrm{\Phi }/f),V^\mu =(\xi ^\mu \xi ^{}+\xi ^{}^\mu \xi )/2,A^\mu =(\xi ^\mu \xi ^{}\xi ^{}^\mu \xi )/(2i).`$ The $`3\times 3`$ matrix $`B`$ is the ground state baryon octet, $`m_\mathrm{B}`$ is the common baryon octet mass in the chiral limit and $`\mathrm{\Phi }`$ collects the pseudoscalar meson octet. The next-to-leading-order terms in Eq. (4) have been developed for heavy baryons by Jenkins and Manohar chiral1 . The heavy baryon chiral theory is similar to the non-relativistic formulation of baryon chiral perturbation theory chiral2 . However, the heavy baryon theory has the advantage of manifest Lorentz invariance, and quantum corrections can be computed in a straightforward manner by the ordinary Feynman graphs, rather than the time ordered perturbation theory chiral11 . The Lagrangian has been shown to be suitable for describing the chiral properties of nuclear system in Ref. chiral3 , where one can also find detailed discussions on how to systematically compute the higher order terms of this Lagrangian. In this paper, we limit our calculations up to the squared characteristic small momentum scale $`Q^2`$ (involving no loops) for s-wave $`\eta \text{N}`$ scattering, because the corrections from the higher-order coupling are suppressed, at low energy, by powers of $`Q/\mathrm{\Lambda }_\chi `$ with $`\mathrm{\Lambda }_\chi 1`$ GeV being the chiral symmetry breaking scale. Hence no loops need to be calculated in this paper. If the loop corrections are included, the higher order terms, i.e., next-to-next-to-leading order, should be added. We will consider it in our later work. Expanding $`\mathrm{\Sigma }`$ up to the order of $`1/f^2`$, and using the heavy-baryon approximation, i.e., $$v=\frac{p}{m}=(\sqrt{1+\frac{\text{p}^2}{m^2}},v_x,v_y,v_z)(1,0,0,0)$$ (5) (because $`v_x`$, $`v_y`$ and $`v_z`$ are very small), we easily obtain the Lagrangian for $`\eta `$N interactions: $`_\eta `$ $`=`$ $`{\displaystyle \frac{1}{2}}^\mu \eta _\mu \eta {\displaystyle \frac{1}{2}}\left(m_\eta ^2{\displaystyle \frac{\mathrm{\Sigma }_{\eta \mathrm{N}}}{f^2}}\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}\right)\eta ^2`$ (6) $`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\kappa }{f^2}}\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}^\mu \eta _\mu \eta ,`$ where $`m_\eta `$ corresponds to the mass of $`\eta `$-meson, which is determined by $`m_\eta ^2=\frac{2}{3}B_0(m_q+2m_s)`$. $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$ is the $`\eta \text{N}`$ sigma term, which is determined by $`\mathrm{\Sigma }_{\eta \mathrm{N}}={\displaystyle \frac{2}{3}}[a_1m_q+4a_2m_s+2a_3(m_q+2m_s)].`$ (7) From Eq.(6), we can see that the last three terms of Eq.(3) do not contribute to the $`\eta `$N interactions. The $`\mathrm{\Sigma }_{\eta \mathrm{N}}/f^2`$ term in Eq.(6) is deduced from the first three terms of Eq.(4), which corresponds to the chiral breaking and shifts the effective mass of $`\eta `$-meson in the nuclear medium. The last term of Eq.(6) is the contribution from the last ten terms of Eq.(4), sometimes, which is called “off-shell” term. $`\kappa `$ is a constant relevant to $`d_i`$’s ($`i=1\text{}10`$). Its value is to be determined from the $`\eta `$N scattering length. ### II.2 The determination of the $`\eta `$N sigma term and $`\kappa `$ To calculate $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$, we should know the parameters on the right hand side of Eq. (7). In fact these parameters have been previously discussed in Refs. xis ; xis1 ; xis2 ; xis3 and are used in Ref. chiral01 . As is well known, the KN sigma term can be written as CHL $$\mathrm{\Sigma }_{\mathrm{KN}}=(m_s+m_q)(a_1+2a_2+4a_3)/2.$$ (8) Solving $`a_3`$ from this equation, and then substituting the corresponding expression into Eq. (7) leads to $$\mathrm{\Sigma }_{\eta \mathrm{N}}=\frac{2}{3}\left[\frac{2+r}{1+r}\mathrm{\Sigma }_{\mathrm{KN}}+a_1m_s\left(1\frac{r}{2}\right)a_2m_s(2r)\right],$$ (9) where $`r=m_q/m_s1`$. Expanding the right hand side of Eq. (9) to a Taylor series with respect to $`r`$, we have $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$ $`=`$ $`{\displaystyle \frac{2}{3}}\left(2\mathrm{\Sigma }_{\mathrm{KN}}+a_1m_s2a_2m_s\right)`$ (10) $`{\displaystyle \frac{1}{3}}\left(2\mathrm{\Sigma }_{\mathrm{KN}}+a_1m_s2a_2m_s\right)r`$ $`+\text{higher-order terms in}r.`$ Because of the extreme smallness of $`r`$, and also due to the fact that our formulas are valid merely up to the next-to-leading order, we take only the first two terms, i.e, $`\mathrm{\Sigma }_{\eta \mathrm{N}}=(1/3)(2\mathrm{\Sigma }_{\mathrm{KN}}+a_1m_s2a_2m_s)(2r).`$ Usually, $`r`$ is in the range of (1/24, 1/26) r0 ; r1 ; r2 ; r3 , and we use the modest value $`r=1/25`$. In fact, the concrete value does not matter significantly due to the extreme smallness of $`r`$. The values for $`a_1m_s`$ and $`a_2m_s`$ can be well determined by Gell-Mann Okubo mass formulas, giving the result $`a_1m_s=67`$ MeV. For $`a_2m_s`$, one has 125 MeV chiral or a little bigger value 134 MeV xis3 , and we take the average $`a_2m_s=130`$ MeV. The value for KN sigma term has some uncertainties. The latest result is $`\mathrm{\Sigma }_{\mathrm{KN}}=312\pm 37`$ MeV in the perturbative chiral quark modelKN . The lattice gauge simulation gave $`\mathrm{\Sigma }_{\mathrm{KN}}=450\pm 30`$ MeVxis5 . The result of lattice QCD is $`\mathrm{\Sigma }_{\mathrm{KN}}=362\pm 13`$ MeV KN1 , and prediction of the Nambu-Jona-Lasinio model is $`\mathrm{\Sigma }_{\mathrm{KN}}=425`$ (with an error bar of $`1015`$%)KN2 . Thus, in our calculations, we use $`\mathrm{\Sigma }_{\mathrm{KN}}=380\pm 100`$ MeV in its possible range. Equipped with the above parameters, we finally obtain $`\mathrm{\Sigma }_{\eta \mathrm{N}}=283\pm 131\text{MeV}`$, where $`\pm 131`$ MeV reflects the uncertainty $`\pm 100`$ MeV in $`\mathrm{\Sigma }_{\mathrm{KN}}`$. Naturally, if one uses a smaller $`\mathrm{\Sigma }_{\mathrm{KN}}`$ vakue, e.g., $`\mathrm{\Sigma }_{\mathrm{KN}}=2m_\pi `$ chiral , $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$ would also become smaller. For the other parameter $`\kappa `$, it is not too difficult, from the Lagrangian in Eq. (6), to derive the $`\eta `$N scattering length (on-shell constraints): $`a^{\eta \mathrm{N}}={\displaystyle \frac{1}{4\pi f^2(1+m_\eta /M_\mathrm{N})}}\left(\mathrm{\Sigma }_{\eta \mathrm{N}}+\kappa m_\eta ^2\right).`$ (11) So we can determine $`\kappa `$ with a given $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$ and $`a^{\eta \mathrm{N}}`$ via the relation $`\kappa =4\pi f^2\left({\displaystyle \frac{1}{m_\eta ^2}}+{\displaystyle \frac{1}{m_\eta M_\mathrm{N}}}\right)a^{\eta \mathrm{N}}{\displaystyle \frac{\mathrm{\Sigma }_{\eta \mathrm{N}}}{m_\eta ^2}}.`$ (12) Recently, Green et al. sc analyzed the new experimental data from GRAAL gall , and gave the real part of $`\eta \text{N}`$ scattering length $`a^{\eta \mathrm{N}}=0.91`$ fm, which agrees to their previous result sc0 . With the similar method, Arndt et alttt also predicted $`a^{\eta \mathrm{N}}=1.031.14`$ fm, comparable to that found by Green et al. So one can assume that $`a^{\eta \mathrm{N}}`$ is in the range of 0.91 $``$ 1.14 fm. Using the central value $`a^{\eta \mathrm{N}}=1.02\text{fm}`$ leads to $`\kappa =0.40\pm 0.08`$ fm. For the eta and nucleon masses, we use $`m_\eta =547.311`$ MeV mass and $`M_\mathrm{N}=939`$ MeV. It should be pointed out that the $`\eta `$N interactions in the present model come from the term of $`\mathrm{\Sigma }_{\eta N}/f^2`$ and the off-shell term, while the leading Tomozawa-Weinberg term simply vanishes. We do not consider any other non-diagonal coupled channel, which was investigated with the chiral coupled channel model by Waas and Weise b3 . According to their calculations, the contribution of non-diagonal coupled channel to the $`\eta \mathrm{N}`$ optical potential is on the order of $`20`$ MeV at normal nuclear density. ## III In-medium properties of $`\eta `$-meson The Lagrangian for one $`\eta `$-meson in nuclear matter is given by $`=_0+_\eta ,`$ (13) where $`_0`$ is the Lagrangian for the nucleon system. In this paper, we adopt the standard Lagrangian, $`_0`$, for the nucleon system in relativistic mean-field theory (given in the appendix). $`_\eta `$ is the Lagrangian for $`\eta `$-meson, which is given in Eq. (6). On application of the Lagrangian in Eq. (13), we immediately have the equation of motion for the $`\eta `$-meson field $`\left(_\mu ^\mu +m_\eta ^2{\displaystyle \frac{\mathrm{\Sigma }_{\eta N}}{f^2}}\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}+{\displaystyle \frac{\kappa }{f^2}}\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}_\mu ^\mu \right)\eta =0.`$ (14) Defining the $`\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}`$ fluctuation $`\delta `$ as $`\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}=\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}+\delta ,`$ (15) where $`\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}`$ is the vacuum expectation value. Because the mean-field approximation is a very familiar method which has already been used in studying the in-medium properties of kaons with a similar chiral approach sc3 ; ch , We adopt it in our present calculations. At the mean-field level, we neglect the fluctuation $`\delta `$. Then the equation of motion for the $`\eta `$-meson field is simplified to $$\left(_\mu ^\mu +m_\eta ^2\frac{\mathrm{\Sigma }_{\eta \mathrm{N}}}{f^2}\rho _s+\frac{\kappa }{f^2}\rho _s_\mu ^\mu \right)\eta =0,$$ (16) where $`\rho _s\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}`$ is the scalar density. Plane wave decomposition of Eq. (16) yields $`\omega ^2+\stackrel{}{k}^2+m_\eta ^2{\displaystyle \frac{\mathrm{\Sigma }_{\eta \mathrm{N}}}{f^2}}\rho _s+{\displaystyle \frac{\kappa }{f^2}}\rho _s\left(\omega ^2+\stackrel{}{k}^2\right)=0.`$ (17) The $`\eta `$-meson effective mass, $`m_\eta ^{}`$, in the nuclear medium is defined by $$\omega =\sqrt{m_{\eta }^{}{}_{}{}^{2}+\stackrel{}{k}^2}.$$ (18) Substituting this equation into Eq. (17) leads to an explicit expression $`m_\eta ^{}=\sqrt{\left(m_\eta ^2{\displaystyle \frac{\mathrm{\Sigma }_{\eta \mathrm{N}}}{f^2}}\rho _s\right)/\left(1+{\displaystyle \frac{\kappa }{f^2}}\rho _s\right)}.`$ (19) Simultaneously, the last two terms on the right hand side of Eq. (17) is the $`\eta `$-meson self-energy, i.e., $`\mathrm{\Pi }(\omega ,\stackrel{}{k};\rho _s)={\displaystyle \frac{\mathrm{\Sigma }_{\eta \mathrm{N}}}{f^2}}\rho _s+{\displaystyle \frac{\kappa }{f^2}}\rho _s\left(\omega ^2+\stackrel{}{k}^2\right),`$ (20) which is a function of the $`\eta `$-meson single-particle energy $`\omega `$ and the momentum $`\stackrel{}{k}`$. Accordingly, the optical potential for $`\eta `$-meson in the nuclear matter is given by $$U_\eta =\frac{1}{2m_\eta }\mathrm{\Pi }(\omega ,\stackrel{}{k}=0;\rho _s)=\frac{m_{\eta }^{}{}_{}{}^{2}m_\eta ^2}{2m_\eta }.$$ (21) To obtain the $`\eta `$-meson in-medium properties, we need a relation between the scalar density $`\rho _s`$ and the nucleon density $`\rho _\mathrm{N}=\mathrm{\Psi }_\mathrm{N}^{}\mathrm{\Psi }_\mathrm{N}`$. Because there is only one single $`\eta `$-meson in the nuclear matter, its effect on the nuclear matter is neglectable. According to the relativistic mean-field theory, we have the following relation between $`\rho _s`$, $`\rho _\mathrm{N}`$, and the $`\sigma `$ mean-field value $`\sigma _0`$: $$\rho _s=\left(M_\mathrm{N}+g_\sigma ^\mathrm{N}\sigma _0\right)^3f(x),$$ (22) where the function $`f(x)`$ is defined to be $$f(x)\left[x\sqrt{1+x^2}\mathrm{ln}\left(1+\sqrt{1+x^2}\right)\right]/\pi ^2$$ (23) with $`x`$ being the ratio of the nucleon’s Fermi momentum to its effective mass, i.e., $$x\frac{k_\mathrm{F}}{M_\mathrm{N}^{}}=\left(\frac{3}{2}\pi ^2\rho _\mathrm{N}\right)^{1/3}/\left(M_\mathrm{N}+g_\sigma ^\mathrm{N}\sigma _0\right).$$ (24) The mean-field value $`\sigma _0`$ is connected to the scalar density $`\rho _s`$ by $$\rho _s=\left(m_\sigma ^2\sigma _0+g_2\sigma _0^2+g_3\sigma _0^3\right)/g_\sigma ^\mathrm{N}.$$ (25) Therefore, for a given nucleon density $`\rho `$, we can first solve $`\sigma _0`$ from $$m_\sigma ^2\sigma _0+g_2\sigma _0^2+g_3\sigma _0^3=g_\sigma ^\mathrm{N}\left(M_\mathrm{N}+g_\sigma ^\mathrm{N}\sigma _0\right)^3f(x),$$ (26) and then calculate the scalar density $`\rho _s`$ from Eq. (25) or (22). The detailed derivation of the Eqs. (22)–(26) can be seen in Ref. rho . To be self-contained, we also attach a brief derivation in the appendix. In numerical calculations, we adopt the NL3 parameter set Lalazissis97prc55 i.e., $`m_\sigma =508.194`$ MeV, $`m_\omega =782.501`$ MeV, $`g_\sigma ^\mathrm{N}=10.217`$, $`g_\omega ^\mathrm{N}=12.868`$, $`g_2=10.434\mathrm{fm}^1`$ and $`g_3=28.885`$. The numerical results for $`\rho _s`$-$`\rho _\mathrm{N}`$ are given in Fig. 1, where one can see clearly that $`\rho _s`$ is an increasing function of the nuclear density. When the density is about 1.5 times lower than the nuclear saturation density, $`\rho _s`$ is nearly proportional to $`\rho _\mathrm{N}`$. However, when the density is about 2 times higher than the normal nuclear density, $`\rho _s`$ is nearly a constant. The mean-field value of the sigma filed is also given in Fig. 1 with a dotted curve. Its density behaviour is similar to that of $`\rho _s`$. ## IV Results and discussions In this section, we discuss the effective mass, optical potential in nuclear medium, and the off-shell behavior of $`\eta `$-meson, respectively. For zero momentum $`\eta `$-meson, we can see, from Eq. (18), that the energy $`\omega `$ is equal to its effective mass. Therefore, we do not mention the $`\eta `$-meson energies any more in the following discussions. In the calculation, the precision of the $`\eta `$-meson effective mass and optical potential are determined by the two parameters $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$ and $`\kappa `$. Equation (12) connects the parameter $`\kappa `$ to the scattering length $`a^{\eta \mathrm{N}}`$, whose possible values are collected in Tab. 1. To reflect uncertainties in the two quantities $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$ and $`a^{\eta \mathrm{N}}`$, we take the sigma term $`\mathrm{\Sigma }_{\eta N}=`$ 150, 280, and 410 MeV, and the scattering length $`a^{\eta \mathrm{N}}=0.91`$ sc and 1.04 ttt fm, in numerical calculations. Figs. 2 and 3 show the $`\eta `$-meson effective mass and nuclear optical potential of $`\eta `$-meson as functions of the nuclear density. The results from Ref. b3 (strait line) is also shown in Fig. 2 for comparison. The curves in Figs. 2 and 3 are obviously divided into three groups which correspond to different scattering lengths $`a^{\eta \mathrm{N}}=`$ 0.91, 1.14 fm and $`\kappa =0`$, respectively. The dotted, solid and dash-dotted curves in each group correspond to $`\mathrm{\Sigma }_{\eta \mathrm{N}}=`$ 150, 280, and 410 MeV, respectively. ### IV.1 Effective mass It is obvious, from Fig. 2, that the $`\eta `$-meson effective mass decreases almost linearly in the region $`\rho <\rho _0`$. In this region, the results of Ref. b3 also show a linear relation for the effective masse with nuclear density. At higher densities, however, the effective mass decreases non-linearly, and the decreasing speed becomes smaller and smaller and at last nearly constant in the range $`\rho >2\rho _0`$. The reason is that, when the density is higher than about 2 times the normal nuclear saturation density, $`\rho _s`$ nearly is a constant (see Fig. 1). For the same scattering length, we find that, at low density region $`\rho 0.5\rho _0`$, the effective mass is nearly independent of the sigma term $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$. When we set $`\mathrm{\Sigma }_{\eta \mathrm{N}}=280\pm 130`$ MeV, which changes in a large range, the variation of the effective mass is within $`\pm 4`$ MeV at $`\rho =\rho _0`$. And at high nuclear density, say $`\rho =3\rho _0`$, the variation is within $`\pm 10`$ MeV compared with that at the central value of $`\mathrm{\Sigma }_{\eta N}=280`$ MeV. Thus, we can conclude that the effective mass of $`\eta `$-mesons is insensitive to the concrete value of $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$ at low density region. Although the latest predictions sc ; ttt give large scattering lengths $`a^{\eta N}=0.911.14`$ fm, there are other different predictions s25 ; s27 ; s30 ; sl46 ; s48 ; s51 ; s55 ; s62 ; s68 ; s71 ; s75 ; sd75 ; sc0 ; s91 ; s98 ; s99 . To see the effects of different scattering length values on the $`\eta `$ effective mass, we show, in Fig. 2, the results for $`a^{\eta N}=0.91`$ and 1.14 fm, respectively. On the other hand, in Tab. I, we give, at normal nuclear density, the effective mass corresponding to the respective $`\eta `$N scattering length in the literature s25 ; s27 ; s30 ; sl46 ; s48 ; s51 ; s55 ; s62 ; s68 ; s71 ; s75 ; sd75 ; s91 ; s98 ; s99 ; sc0 ; sc ; ttt . From Fig. 2, we find that, with the same sigma term $`\mathrm{\Sigma }_{\eta N}`$, the effective mass depends strongly on the scattering length $`a^{\eta N}`$. At $`\rho =\rho _0`$, the effective mass (with $`\mathrm{\Sigma }_{\eta N}=280`$MeV) is $`m_\eta ^{}/m_\eta =0.85`$ for $`a^{\eta N}=0.91`$ and $`m_\eta ^{}/m_\eta =0.825`$ for $`a^{\eta N}=1.14`$ fm. When the scattering length varies from 0.25 fm to 1.14 fm, the effective mass will run from 0.95$`m_\eta `$ to 0.825$`m_\eta `$. Corresponding to $`a^{\eta N}=0.911.14`$ fm, which are favored by recent works, and $`\mathrm{\Sigma }_{\eta \mathrm{N}}`$, which is predicted in Sec. II, the effective mass is $`(0.84\pm 0.015)m_\eta `$. At normal nuclear density, the effective mass in Ref. b3 is 0.95$`m_\eta `$, which agrees with result of the small scattering length $`a^{\eta N}=0.25`$ fm. As pointed out in the above, the effective mass changes nonlinearly with increasing densities in the region $`\rho _0<\rho <2\rho _0`$. This behavior agrees to the predictions by Tsushima et al. efm with quark-meson coupling model. The effective mass at $`\rho =\rho _0`$ predicted by them is about 0.88$`m_\eta `$, which just corresponds to the result with scattering length $`a^{\eta N}=0.68`$ fm. This can be clearly seen from Tab. I . The outstanding characteristic of our results is that the present calculations give much smaller effective mass than the others when we adopt the larger scattering length. It should be mentioned that the chiral coupled channel model b3 gives much larger in-medium effective mass for $`\eta `$-mesons than our predictions. The main reason is as such. In the chiral coupled channel model, there are only the leading-order terms, and so, the contributions to the effective mass come only from the non-diagonal coupled channel. While in our model, the leading-order terms do not contribute to the calculations. All the contributions to the results come from the next-to-leading-order terms. ### IV.2 Optical potential The optical potential $`U_\eta `$ as a function of nuclear density is plotted in Fig. 3. We find that the density behavior of $`U_\eta `$ is quite similar to the effective mass in Fig. 3. The reason is that the optical potential has a relation $`U_\eta m_\eta ^{}m_\eta `$ as an approximation, which varies linearly with the effective mass $`m_\eta ^{}`$ of $`\eta `$-meson. Similarly, it is also seen that the effect from the uncertainties of sigma term $`\mathrm{\Sigma }_{\eta N}`$ are quite limited in its possible range, and the optical potential depend strongly on the value of the scattering length. At normal nuclear density, the upper limit of the uncertainties from the sigma term $`\mathrm{\Sigma }_{\eta N}`$ is no more than $`8`$ MeV. However, the optical potential can change from $`78`$ MeV to $`88`$MeV, when we modify the scattering length $`a^{\eta N}`$ from 0.91 to 1.14 fm. Because there are still uncertainties for the $`\eta `$N scattering length, we listed the possible potential depths corresponding to the possible scattering lengths appearing in literature in Tab. I. From the table, we can see that the potential depth at normal nuclear density ranges from 26 MeV to 88 MeV, because of the uncertainties of scattering lengths. According to the newest predictions, i.e. $`a^{\eta N}=0.911.14`$ fmsc ; ttt , the potential depth is about $`83\pm 5`$ MeV. This is a very strong attractive potential which was never predicted by the previous models. There have been some predictions for the nuclear potential of $`\eta `$-mesons in other references. According to the SU(3) chiral dynamics with coupled channels, the optical potential depth at normal nuclear density is $`U_\eta 20`$ MeV b3 , which is close to our formulas with a smaller scattering length $`a_{\eta N}<0.25`$ fm. In Ref. op , by assuming that the mass of the $`N^{}(1535)`$ did not change in the medium, the optical potential $`U_\eta =34`$ MeV was obtained, which is close to our calculation with $`a_{\eta N}0.30`$ fm. The $`\eta `$ potential from QMC model by Tsushima et al. and chiral unitary approach by Inoue et al., are typically $`60`$ MeV and $`54`$ MeV, which are comparable to our formulas with $`a_{\eta N}=0.550.68`$ fm. Therefore, if we want to obtain shallower optical potential, we need to use a smaller scattering length. Because recent works favor the bigger scattering length, our formulas give much deeper optical potential. ### IV.3 The effect of off-shell term Finally, we discuss the role of the off-shell term in our calculation. In present model, the off-shell term $`\kappa `$ is determined by the scattering length $`a^{\eta N}`$. From the analysis of the subsections A and B, we know that the scattering length $`a^{\eta \mathrm{N}}`$ strongly affects the calculations. The importance of the off-shell behavior for low energy scattering had been pointed out in many Refs.off ; chiral ; sc3 . To clarify the effect of off-shell term on our calculation thoroughly, we turn off the off-shell term ($`\kappa =0`$) and show the results in Fig. 2 and 3. At $`\rho =\rho _0`$, without the off-shell terms, the effective mass is $`m_\eta ^{}/m_\eta 0.94\pm 0.03`$, and the optical potential is $`(32\pm 16)`$ MeV corresponding to $`\mathrm{\Sigma }_{\eta N}=280\pm 130`$ MeV. Thus, without the off-shell terms, we no longer have strong attractive potential for $`\eta `$-meson in nuclear medium. The calculations are independent of the scattering length. Also in this case, the calculations depend strongly on the quantity of $`\mathrm{\Sigma }_{\eta N}`$. Without the off-shell terms, the variation of the optical potential from the uncertainties of $`\mathrm{\Sigma }_{\eta N}`$ can reach about 30 MeV at normal nuclear density. However, it is no more than 8 MeV, when the off-shell behavior is considered. Thus, the off-shell terms can dramatically depress the effects from the uncertainties of $`\mathrm{\Sigma }_{\eta N}`$. ## V Summary In this paper, we have derived an effective Lagrangian for $`\eta `$N s-wave interaction from the effective meson-baryon chiral Lagrangian including the next-to-leading-order terms. Up to $`1/f^2`$ terms for s-wave $`\eta `$N interaction, only the sigma term and off-shell term survive. It is found that the $`\eta `$N sigma term is $`\mathrm{\Sigma }_{\eta N}=280\pm 130`$ MeV according to the KN sigma term. The off-shell term $`\kappa `$ is determined by the scattering length. If we adopt the newest predictions $`a_{\eta N}0.911.14`$ fm for the scattering lengthssc ; ttt , we obtain the value $`\kappa =0.40\pm 0.08`$ fm. Combining the relativistic mean-field theory for nucleon system, we calculate the effective mass and optical potential of $`\eta `$-mesons in uniform nuclear medium in the mean-field approximation. According to the latest predictions $`a_{\eta N}0.911.14`$ fm for the scattering lengthssc ; ttt , at normal nuclear density the effective mass is about $`(0.84\pm 0.015)m_\eta `$, and the depth of optical potential is $`U_\eta (83\pm 5)`$ MeV. Finally, we should point out the importance of the next-to-leading-order terms of the chiral Lagrangian again. In fact, the leading-order terms do not contribute to the $`\eta `$N interactions. All contribution comes from the next-to-leading order terms. It indicates that the next-to-leading order terms should be included in the study of the $`\eta `$N interaction. In the present paper, we do not consider corrections from the non-diagonal coupled channel. According the study of Waas and Weise, the correction may be on the order of 20 MeV for the optical potential. ## Acknowledgements This work was supported by the Natural Science Foundation of China (10275037, 10375074, 10575054, 90203004) and the Doctoral Programme Foundation of the China Institution of Higher Education (20010055012). We thank Prof. X. B. Zhang for helpful discussions. ## Appendix A relation between the scalar density and nucleon density in the relativistic mean-field approach In this appendix, we give a short derivation of the relation between the scalar density and nucleon density in the relativistic mean-field approach (RMF). In RMF, the effective Lagrangian density rmf can be written as $`_0`$ $`=`$ $`\overline{\mathrm{\Psi }}_\mathrm{N}(i\gamma ^\mu _\mu M_\mathrm{N})\mathrm{\Psi }_\mathrm{N}g_\sigma ^N\overline{\mathrm{\Psi }}_\mathrm{N}\sigma \mathrm{\Psi }_\mathrm{N}`$ (27) $`g_\omega ^N\overline{\mathrm{\Psi }}_\mathrm{N}\gamma ^\mu \omega _\mu \mathrm{\Psi }_\mathrm{N}g_\rho ^N\overline{\mathrm{\Psi }}_\mathrm{N}\gamma ^\mu \rho _\mu ^a{\displaystyle \frac{\tau _a}{2}}\mathrm{\Psi }_\mathrm{N}`$ $`+{\displaystyle \frac{1}{2}}^\mu \sigma _\mu \sigma {\displaystyle \frac{1}{2}}m_\sigma ^2\sigma ^2{\displaystyle \frac{1}{3}}g_2^2\sigma ^3{\displaystyle \frac{1}{4}}g_3^2\sigma ^4`$ $`{\displaystyle \frac{1}{4}}\mathrm{\Omega }^{\mu \nu }\mathrm{\Omega }_{\mu \nu }+{\displaystyle \frac{1}{2}}m_\omega ^2\omega ^\mu \omega _\mu {\displaystyle \frac{1}{4}}R^{a\mu \nu }R_{\mu \nu }^a`$ $`+{\displaystyle \frac{1}{2}}m_\rho ^2\rho ^{a\mu }\rho _\mu ^a{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }`$ $`e\overline{\mathrm{\Psi }}_\mathrm{N}\gamma ^\mu A^\mu {\displaystyle \frac{1}{2}}(1+\tau _3)\mathrm{\Psi }_\mathrm{N}`$ with $`\mathrm{\Omega }^{\mu \nu }=^\mu \omega ^\nu ^\nu \omega ^\mu ,R^{a\mu \nu }=^\mu \rho ^{a\nu }^\nu \rho ^{a\mu },F^{\mu \nu }=^\mu A^\nu ^\nu A^\mu .`$ On application of the mean-field approximation, we have the equation of motion for nucleons: $$\left(\gamma _\mu k^\mu M_\mathrm{N}g_\sigma ^\mathrm{N}\sigma _0g_\omega ^\mathrm{N}\gamma ^0\omega _0g_\rho ^\mathrm{N}\gamma ^0\tau ^3\rho _{0}^{}{}_{3}{}^{}\right)\mathrm{\Psi }_\mathrm{N}=0.$$ (28) where the $`\sigma `$, $`\omega `$, and $`\rho `$ fields are replaced with their mean-field values $`\sigma _0`$, $`\omega _0`$ and $`\rho _0`$. $`\sigma _0`$ and $`\omega _0`$ satisfy $`m_\sigma ^2\sigma _0+g_2\sigma _0^2+g_3\sigma _0^3`$ $`=`$ $`g_\sigma ^\mathrm{N}\rho _s,`$ (29) $`m_\omega ^2\omega _0`$ $`=`$ $`g_\omega ^\mathrm{N}\rho _\mathrm{N},`$ (30) with $`\rho _s\overline{\mathrm{\Psi }}_\mathrm{N}\mathrm{\Psi }_\mathrm{N}`$ and $`\rho _\mathrm{N}\mathrm{\Psi }_\mathrm{N}^{}\mathrm{\Psi }_\mathrm{N}.`$ Therefore, at the mean-field level, the energy density of nuclear matter is $`\epsilon `$ $`=`$ $`{\displaystyle \frac{1}{2}}m_\sigma ^2\sigma _0^2+{\displaystyle \frac{1}{3}}g_2\sigma _0^3+{\displaystyle \frac{1}{4}}g_3\sigma _0^4+{\displaystyle \frac{1}{2}}m_\omega ^2\omega _0^2`$ (31) $`+{\displaystyle \frac{4}{(2\pi )^3}}{\displaystyle _0^{k_\mathrm{F}}}\left(\stackrel{}{k}^2+M_{\mathrm{N}}^{}{}_{}{}^{2}\right)^{1/2}d\stackrel{}{k},`$ where $`M_\mathrm{N}^{}=M_\mathrm{N}+g_\sigma ^\mathrm{N}\sigma _0`$ is the effective mass of nucleons. In Eq. (31), the energy density has been expressed as an explicit function of $`\sigma _0`$. Because $`\sigma _0`$ should minimize $`\epsilon `$, i.e., $`\epsilon (\sigma _0)/\sigma _0`$, we immediately have $$m_\sigma ^2\sigma _0+g_2\sigma _0^2+g_3\sigma _0^3=\frac{4g_\sigma ^\mathrm{N}}{(2\pi )^3}_0^{k_\mathrm{F}}\frac{M_\mathrm{N}^{}\mathrm{d}\stackrel{}{k}}{(\stackrel{}{k}^2+M_{\mathrm{N}}^{}{}_{}{}^{2})^{1/2}},$$ (32) which is nothing but the Eq. (26). Equation (29) corresponds to Eq. (25). And comparing Eq. (32) with Eq. (29) then gives the Eq. (22).
warning/0506/astro-ph0506092.html	ar5iv	text	# QPOs and firehose instabilities in neutron star magnetospheres in accreting systems Also at \]Institute for Advanced Studies in Basic Sciences, Zanjan 45195, Iran. ## Abstract We show that the interaction of an accretion disk with the magnetosphere of a neutron star can excite resonant shear Alfvén waves with Hz-kHz frequencies in a region of enhanced density gradients. This is the the region where accretion material flows along the magnetic field lines in the magnetosphere. We argue that due to the pressure anisotropy produced by the plasma flow, firehose instabilities are likely to occur. Furthermore, for a dipolar field topology, we show that a new instability develops due to both magnetic field curvature and the plasma flow. Recent observations of quasi-periodic oscillations (QPOs), particularly kHz QPOs, in X-ray emissions from accreting binaries, have aroused a lot of interest in the astrophysical community. The QPOs are very strong and remarkably coherent with frequencies ranging from $`10`$ Hz to $`1200`$ Hz. They have been observed in the X-ray flux of about 20 accreting neutron star sources and five black hole sources by Rossi X-Ray Timing Explorer. Almost all sources have shown twin spectral peaks in the QPOs in the kHz part of the X-ray spectrum, with the value of the peak separation being anti-correlated with the QPO frequencies van der Klis (2000); Di Salvo, Méndez & van der Klis (2003); Migliari, van der Klis & Fender (2003). The clear similarities of kHz QPO properties in black hole systems to those in neutron star binaries Psaltis, Belloni, & van der Klis (1999); Belloni, Psaltis, & van der Klis (2002), and in white dwarf systems Warner, Woudt, & Pretorius (2003), suggest that the oscillation likely originate in the accretion flow surrounding the central object. Motivated by this argument, a variety of accretion-based models have been proposed. Most models, however, fail to provide a general explanation of QPO features in all potential sources. See, for example, Li & Narayan Li & Narayan (2004) and references therein. Li & Narayan Li & Narayan (2004) studied Rayleigh-Taylor and Kelvin-Helmholtz instabilities at a possible interface between the star’s magnetosphere and the accretion disk. They found that modes with low order azimuthal wavenumbers are expected to grow to large amplitude and to contribute to kHz QPOs. In their study, the magnetic field allows an interface with abrupt spatial discontinuities in the flow density and/or angular velocity across the interface. They ignore, however, the dynamical role of the magnetic field of the central object and its interaction with the accretion disk/flow. Based on theoretical models and observations of the aurora in the Earth’s magnetosphere, Rezania et al. Rezania et al. (2004) and Rezania & Samson Rezania & Samson (2005) (hereafter paper I) have recently proposed a generic magnetospheric model for accretion disk-neutron star systems to address the occurrence and the behavior of the observed QPOs in those systems. In the Earth’s magnetosphere, the occurrence of aurora is a result of the resonant coupling between shear Alfvén waves and fast compressional waves (produced by the solar wind). These resonances are known as field line resonances (FLRs). Paper I argued that this resonant coupling is also likely to occur in neutron star magnetospheres, due to interaction with the accreting plasma. The MHD interaction of the infalling plasma (with a sonic/supersonic speed) with the neutron star magnetosphere, would alter not only the plasma flow toward the surface of the star, as assumed by current QPO models, but also the structure of the star’s magnetosphere. The magnetic field of the neutron star is distorted inward by the infalling plasma of the Keplerian accretion flow. Furthermore, the plasma would likely be able to penetrate through magnetic field lines and produce enhanced density regions within the magnetosphere (similar to the magnetospheric interface introduced in Li & Narayan (2004)). Any instability in the compressional action of the accretion flow would alter the quasi-equilibrium pressure balance between the inward pressure of the infalling flow $`\rho v_r^2/2`$ and the outward magnetic pressure $`B_p^2/8\pi `$. This process then can excite perturbations in the enhanced density region. Here $`\rho `$ and $`v_r`$ are the density and radial velocity of the infalling matter and $`B_p`$ is the poloidal magnetic field on the plane of the disk. In analogy with the interaction of the solar wind with the earth’s magnetosphere Samson (1991), one would expect the excitation of resonant shear Alfvén waves, or FLRs, Rezania & Samson (2005). Paper I assumed a simple geometry with a rectilinear magnetic field, and showed that, in the presence of a plasma flow,two resonant MHD modes with frequencies in the kHz range can occur within a few stellar radii. The results in paper I gave a reasonable prediction, both quantitative and qualitative, of the kHz oscillations observed in the X-ray fluxes in X-ray binaries. The existence of a non-zero plasma displacement along the magnetic field lines, which oscillates with resonance frequencies and then modulates the flow of the plasma toward the surface of the neutron star, can explain the kHz quasi-periodicity in the observed X-ray flux from the star. See Rezania & Samson (2005) for more details. In this letter we examine the above problem in a more realistic configuration for the star’s magnetosphere. We base are model on the fact that the topology of the magnetospheres of isolated pulsars is likely close to dipolar. In accreting pulsars, however, the geometry of the magnetosphere will be distorted due to the inward flowing plasma, specifically on the plane of the accretion disk. Nevertheless, we approximate the topology of the magnetosphere of an accreting neutron stars with a dipolar geometry in order to study the structure of shear Alfvén waves in the presence of an ambient flow. Furthermore, due to the existence of ambient flow along the magnetic field, the plasma pressure is not expected to be isotropic. This can be understood by noting that the plasma pressure parallel to the field lines can be defined as $$p_{\|\|}mf(u_{\|\|}v_p)^2d^3u,$$ (1) where $`u_{\|\|}=𝐮𝐁/B`$ is the thermal velocity of the plasma particles and $`v_p`$ is the ambient flow velocity along the field lines. Here $`f(𝐫,𝐮,t)`$ is the plasma distribution function satisfying the Vlasov equation Snyder et al. (1997). Similarly, the perpendicular component of the pressure can be calculated from $$p_{}mfu_{}^2d^3u/2.$$ (2) Hence, the scalar pressure of isotropic MHD must be replaced by a diagonal pressure tensor with two components: a parallel component $`p_{\|\|}`$ acting along the field lines and a perpendicular component $`p_{}`$ acting in the perpendicular direction. The latter can be considered to be the ram pressure. Consequently, in the MHD equations, the flow pressure must be written as: $`𝐩=p_{}𝐈+(p_{}p_{\|\|})\mathrm{𝐛𝐛}`$, where $`𝐈`$ is the identity tensor and $`𝐛=𝐁/B`$ is the unit vector along the magnetic field line. The linearized, perturbed magnetohydrodynamic equations with anisotropic pressure in the presence of an ambient flow now have parallel and perpendicular components: $`\rho \left({\displaystyle \frac{\delta 𝐯}{t}}+𝐯\mathbf{}\delta 𝐯+\delta 𝐯\mathbf{}𝐯\right)_{\|\|}=_{\|\|}\delta p_{\|\|}`$ $`+(p_{}p_{\|\|}){\displaystyle \frac{_{\|\|}\delta B_{\|\|}}{B}}`$ (3a) $`\rho \left({\displaystyle \frac{\delta 𝐯}{t}}+𝐯\mathbf{}\delta 𝐯+\delta 𝐯\mathbf{}𝐯\right)_{}=\mathbf{}_{}\delta \left(p_{}+{\displaystyle \frac{B^2}{8\pi }}\right)`$ $`+\mathrm{\Xi }\left({\displaystyle \frac{\delta 𝐁\mathbf{}𝐁+𝐁\delta \mathbf{}𝐁}{4\pi }}\right)_{}`$ (3b) $`{\displaystyle \frac{\delta 𝐁}{t}}=\mathbf{}\times (\delta 𝐯\times 𝐁+𝐯\times \delta 𝐁),`$ (3c) $`\mathbf{}\delta 𝐁=0,`$ (3d) $`\delta p_{\|\|}=2p_{\|\|}\delta B_{\|\|}/B,`$ (3e) $`\delta p_{}=p_{}\delta B_{\|\|}/B,`$ (3f) where $`\delta 𝐯=𝝃/t`$, $`\mathrm{\Xi }=1+2(c_{}^2c_{\|\|}^2)/v_\mathrm{A}^2`$, and we ignore the perturbation in the plasma density, i.e. $`\delta \rho 0`$. Here $`\rho ,p_{\|\|},p_{},𝐯,`$ and $`𝐁`$ are the unperturbed quantities while $`\delta p_{\|\|},\delta p_{},\delta 𝐯,`$ and $`\delta 𝐁`$ are the perturbed quantities. Equations (LABEL:dp\_par) and (LABEL:dp\_per) are calculated from the two equations of state for $`p_{\|\|}`$ and $`p_{}`$ that are known as double adiabatic equations, $`\frac{d}{dt}(p_{\|\|}B^2/\rho ^3)=0`$ and $`\frac{d}{dt}(p_{}/(\rho B))=0`$ Chew et al. (1956). $`\mathrm{\Xi }=1`$ if the pressure is isotropic, i.e $`p_{\|\|}=p_{}=p`$. To avoid complexities, we shall ignore the rotation of the star and consequently neglect both the toroidal field $`B_\varphi `$ and velocity $`v_\varphi `$. These assumptions simplify our calculations significantly, and, we believe, allow the model to retain the important physics. Furthermore, we do not consider a jump condition in the enhanced density region as discussed by Li & Narayan (2004). As a result, we do not address Rayleigh-Taylor and/or Kelvin-Hemholtz instabilities. We expand the MHD equations (QPOs and firehose instabilities in neutron star magnetospheres in accreting systems) in the orthogonal coordinate system ($`\mu ,\nu ,\varphi `$), where $`\mu =\mathrm{cos}\theta /r^2`$ is the magnetic field-aligned coordinate, $`\nu =\mathrm{sin}^2\theta /r`$ numerates magnetic shells in the direction perpendicular to the field line, and $`\varphi `$ the is azimuthal coordinate. The components of the metric are $`h_\mu =h_\nu h_\varphi `$, $`h_\nu =r^2/\mathrm{sin}\theta \sqrt{1+3\mathrm{cos}^2\theta }`$, and $`h_\varphi =r\mathrm{sin}\theta `$ where ($`r,\theta ,\varphi `$) is the spherical coordinate. The metric component $`h_\mu `$ allows a convenient representation of the dipolar magnetic field in the form $`𝐁=\mu ^{\mathrm{mag}}/h_\mu 𝐞_\mu `$ where $`\mu ^{\mathrm{mag}}`$ is the magnetic dipole moment of the star. Note that, for a dipolar magnetic field $`𝐁=B_p𝐞_\mu `$, $`B_ph_\mu =\mu ^{\mathrm{mag}}`$ is constant which leads to $`\mathbf{}\times 𝐁=0`$. We further assume that the ambient flow velocity is along magnetic field lines, ie. $`𝐯=v_p𝐞_\mu `$. Assuming $$\delta (\mu ,\nu ,\varphi ,t)\delta (\nu )e^{ik\mu }e^{i\omega t}e^{im\varphi }.$$ (4) we can reduce the equations of motion (QPOs and firehose instabilities in neutron star magnetospheres in accreting systems) to one second order differential equation for $`\delta v_\nu `$ as $`{\displaystyle \frac{d^2\delta v_\nu }{d\nu ^2}}+F(\mu ,\nu ){\displaystyle \frac{d\delta v_\nu }{d\nu }}+G(\mu ,\nu )\delta v_\nu =0,`$ (5) $`F(\mu ,\nu )={\displaystyle \frac{1}{(c_{}^2+v_A^2)\eta _1}}\left[(c_{}^2+v_A^2)_\nu (\eta _1+\eta _2)+2\mathrm{\Xi }{\displaystyle \frac{v_\mathrm{A}^2}{h_\mu ^2}}\eta _1_\nu \mathrm{ln}h_\mu +{\displaystyle \frac{2ikv_p\eta _1(3c_{\|\|}^2c_{}^2)_\nu \mathrm{ln}h_\mu }{h_\mu (i\omega _D+K_\mu \mathbf{}𝐯_p)}}\right],`$ (5a) $`G(\mu ,\nu )={\displaystyle \frac{1}{(c_{}^2+v_A^2)\eta _1}}[(c_{}^2+v_A^2)_\nu \eta _2+2\mathrm{\Xi }{\displaystyle \frac{v_A^2}{h_\mu ^2}}\eta _2_\nu \mathrm{ln}h_\mu (i\omega _DK_\nu )h_\nu {\displaystyle \frac{h_\nu \mathrm{\Xi }v_A^2}{h_\mu ^2}}{\displaystyle \frac{k^2+(_\mu \mathrm{ln}h_\nu )^2}{i\omega _D+K_\nu \mathbf{}v_p}}`$ $`+{\displaystyle \frac{2v_p_\nu \mathrm{ln}h_\mu /h_\mu }{i\omega _DK_\mu \mathbf{}v_p}}(ik(3c_{\|\|}^2c_{}^2)\eta _2_\nu (h_\mu v_p)/h_\nu )],`$ (5b) $`\eta _1=(h_\varphi Q/h_\mu )/[(i\omega _DK_\mu )Q(m^2/h_\varphi ^2)(c_{}^2+v_\mathrm{A}^2)(i\omega _D+K_\varphi \mathbf{}𝐯_p)],`$ (5c) $`\eta _2=(\eta _1/h_\varphi )\left[{\displaystyle \frac{1}{h_\mu }}_\nu (h_\mu h_\varphi ){\displaystyle \frac{2}{h_\nu }}_\nu h_\mu +{\displaystyle \frac{1}{h_\mu h_\nu }}(h_\mu _\nu v_pv_p_\nu h_\mu ){\displaystyle \frac{ik_\mu h_\nu /h_\nu }{i\omega _D+K_\nu \mathbf{}𝐯_p}}\right],`$ (5d) $`Q=\omega _D^2+K_\varphi ^2+(i\omega _DK_\varphi )\mathbf{}𝐯_p{\displaystyle \frac{v_A^2\mathrm{\Xi }}{h_\mu ^2}}[k^2+{\displaystyle \frac{1}{h_\varphi ^2}}(_\mu h_\varphi )^2],`$ (5e) where $`\omega _D=\omega kv_p/h_\mu `$ is the Doppler shifted frequency, $`K_i=𝐯_p\mathbf{}\mathrm{ln}h_i`$ ($`i=\mu ,\nu ,\varphi `$), $`v_\mathrm{A}=B/\sqrt{4\pi \rho }`$ is the Alfvén wave velocity, $`c_{\|\|}=\sqrt{p_{\|\|}/\rho }`$ and $`c_{}=\sqrt{p_{}/\rho }`$ are sound velocities parallel and perpendicular to the direction of the magnetic field. We note that in deriving Eqs. (5) we assumed that $`_\nu [p_{}+B^2/(8\pi )]0`$. The shear Alfvén resonance happens at $`\eta _1=0`$ or equivalently at $`Q=0`$ leading to $$\omega _D^2+i\mathbf{}𝐯_p\omega _D+K_\varphi ^2K_\varphi \mathbf{}𝐯_p\frac{v_A^2\mathrm{\Xi }}{h_\mu ^2}[k^2+\frac{1}{h_\varphi ^2}(_\mu h_\varphi )^2]=0.$$ (6) As a result, resonance frequencies will be given by $`\omega _\pm =kv_p/h_\mu (i/2)\mathbf{}𝐯_p\pm (1/2)\sqrt{\mathrm{\Delta }},`$ (7) $`\mathrm{\Delta }=4v_A^2\mathrm{\Xi }[k^2+(_\mu \mathrm{ln}h_\varphi )^2]/h_\mu ^2(\mathbf{}𝐯_p2K_\varphi )^2.`$ For a rectilinear configuration, resonance eigenfrequencies Eq. (7) for an incompressible plasma flow along the field lines, $`\mathbf{}𝐯_p=0`$, with an isotropic pressure, $`p_{\|\|}=p_{}=p`$, reduce to ones we obtained in paper I, and as expected: $`\omega _\pm =k(v_p\pm v_A)`$. Equation (7), however, shows that whenever the ambient flow is compressible, i.e. $`\mathbf{}𝐯_p0`$, and/or $`\mathrm{\Delta }<0`$, waves do not propagate and an instability develops. In general, the ambient plasma is fairly incompressible, i.e. $`\mathbf{}𝐯_p=0`$. However, due to the topological deformation of the magnetosphere caused by the compressional action of accreting material, a non-zero density gradient through the plasma, and so $`\mathbf{}𝐯_p0`$ can be expected: $`\mathbf{}𝐯_p={\displaystyle \frac{1}{\rho }}{\displaystyle \frac{d\rho }{dt}}={\displaystyle \frac{1}{\rho }}{\displaystyle \frac{\rho }{t}}{\displaystyle \frac{1}{\rho }}𝐯_p\mathbf{}\rho ,`$ $`𝐯_p\mathbf{}\mathrm{ln}\rho .`$ (8) As a result, the resonant mode will grow (decay) if $`𝐯_p\mathbf{}\rho >0`$ ($`<0`$). Approximating the plasma inflow velocity with the free fall velocity $`v_pv_{\mathrm{ff}}(r)=(2GM/r)^{1/2}`$, the growth/decay timescale will be as order of $`\tau r/v_p=(r^3/2GM)^{1/2}6\times 10^5(M/M_{})^{1/2}(r/10\mathrm{km})^{3/2}`$ s. Therefore, the closer to the star the faster the instability develops. The condition $`\mathrm{\Delta }<0`$ is satisfied whether $`\mathrm{I}:\mathrm{\Xi }<0c_{\|\|}^2>c_{}^2+v_\mathrm{A}^2/2,`$ (9a) $`\mathrm{II}:(\mathbf{}𝐯_p2K_\varphi )^2>4v_A^2\|\mathrm{\Xi }\|[k^2+(_\mu \mathrm{ln}h_\varphi )^2]/h_\mu ^2.`$ Case I, that is known as the firehose instability in the literature, happens when $`p_{\|\|}`$ is much larger than $`p_{}+p_\mathrm{M}`$, where $`p_\mathrm{M}=B^2/(8\pi )`$ is the magnetic pressure. The magnetic field channeling the parallel plasma streams experiences a similar instability. Whenever the flux tube is slightly bent, the flowing plasma exerts a centrifugal force, that tends to enhance the initial bending. The field line bending is proportional to the density of energy in plasma motion along the magnetic field $`\rho v_{\|\|}^2p_{\|\|}`$. Recalling Eq. (1), the ambient plasma flow would enhance the firehose instability in the magnetosphere of an accreting neutron star: $$c_{\|\|}^2c_{}^{}{}_{\|\|}{}^{2}+v_p^2,$$ (10) where $`c_{}^{}{}_{\|\|}{}^{2}(m/\rho )fu_{\|\|}^2d^3u`$ and we assume that a Maxwellian distribution function, so the cross term vanishes. Inserting Eq. (10) into Eq. (9a), we find $$v_p^2>v_\mathrm{A}^2/2+(c_{}^2c_{}^{}{}_{\|\|}{}^{2}).$$ (11) Therefore, a firehose instability develops whenever condition (11) is satisfied. For an isotropic pressure, i.e. $`\mathrm{\Xi }1`$, the firehose instability would not be expected. However, an instability may arise if the condition (LABEL:Delta1) is satisfied. For an incompressible flow, we find that Eq. (LABEL:Delta1) reduces to $$v_p>v_\mathrm{A}\sqrt{1+q^2},$$ (12) where $`q=k/(_\mu \mathrm{ln}h_\varphi )`$. It is necessary to note that this instability is enhanced whenever both $`v_p0`$ and $`_\mu \mathrm{ln}h_\varphi 0`$. The latter condition is due to the curvature of magnetic field lines. Therefore, the non-flat topology of the magnetic field can trigger some MHD instabilities through the magnetosphere. An interesting note is that the condition (LABEL:Delta1) is only valid for for an isotropic pressure flow. When the wave transfers energy to the flow, the wave decays, and the extraction of flow energy by the wave will lead to a growing mode. For a superAlfvénic flow, the extraction of energy from the flow and growing MHD waves is very likely Joarder et al. (1997). In this case also the growth/decay timescale will be as order of $`\tau r/v_p=(r^3/2GM)^{1/2}6\times 10^5(M/M_{})^{1/2}(r/10\mathrm{km})^{3/2}`$ s. Furthermore, by approximating the Alfvén velocity by $`v_\mathrm{A}B(r)/\sqrt{4\pi \rho _{\mathrm{ff}}}`$ where $`\rho _{\mathrm{ff}}=\dot{M}/(v_{\mathrm{ff}}4\pi r^2)`$ is the free fall mass density ($`v_pv_{\mathrm{ff}}(r)=\sqrt{2GM/r}`$), the instability condition (12) is satisfied for distance further than $$r>1.7\times 10^6\mathrm{cm}(1+q^2)^{2/7}\mu _{26}^{4/7}\dot{M}_{17}^{2/7}(M/M_{})^{1/7}.$$ (13) Therefore, this instability is very likely to develop at a position where the accretion disk can distort the dipolar magnetosphere. Here $`\mu _{26}`$ is the magnetic field dipole moment at the surface of star in units of $`10^{26}`$ G cm<sup>3</sup> and $`\dot{M}_{17}`$ is the mass of accretion rate in units of $`10^{17}`$ g s<sup>-1</sup>. We calculate Eq. (13) for a solar mass neutron star with 10 km radius. To summarize, we believe we have shown that there are generic plasma instabilities associated with the radial position where the inflowing material of the accretion disk leads to a distortion of the inner magnetosphere of the neutron star and field aligned plasma flows. We suspect that these instabilities may be relevant in the understanding of some details of observed QPOs in X-ray binary systems, particularly when linked to shear Alfvén waves. The MHD interaction of the infalling plasma with the neutron star’s magnetosphere can alter the topology of the star’s inner magnetoshere. The plasma flows and topology lead to the possibility of : (1) firehose instabilities associated the pressure anisotropy produced by the plasma flow; (2) convective growth of waves in the plasma flow. The unstable modes might produce shear Alfvén resonances with large amplitude giving strong quasi-periodic variations in X-ray fluxes. ###### Acknowledgements. This research was supported by the Natural Sciences and Engineering Research Council of Canada (NSERC). —————————————————————————————-
warning/0506/cond-mat0506053.html	ar5iv	text	# Pairs of SAT Assignment in Random Boolean Formulæ ## 1 Introduction and outline Consider a string of Boolean variables — or equivalently a string of spins — of size $`N`$: $`\stackrel{}{\sigma }=\{\sigma _i\}\{1,1\}^N`$. Call a $`K`$-clause a disjunction binding $`K`$ of these Boolean variables in such a way that one of their $`2^K`$ joint assignments is set to false, and all the others to true. A formula in a conjunctive normal form (CNF) is a conjunction of such clauses. The satisfiability problem is stated as: does there exist a truth assignment $`\stackrel{}{\sigma }`$ that satisfies this formula? A CNF formula is said to be satisfiable (SAT) if this is the case, and unsatisfiable (UNSAT) otherwise. The satisfiability problem is often viewed as the canonical constraint satisfaction problem (CSP). It is the first problem to have been shown NP-complete cook , i.e. at least as hard as any problem for which a solution can be checked in polynomial time. The $`PNP`$ conjecture states that no general polynomial-time algorithm exists that can decide whether a formula is SAT or UNSAT. However formulas which are encountered in practice can often be solved easily. In order to understand properties of some typical families of formulas, one introduces a probability measure on the set of instances. In the random $`K`$-SAT problem, one generates a random $`K`$-CNF formula $`F_K(N,M)`$ as a conjunction of $`M=N\alpha `$ $`K`$-clauses, each of them being uniformly drawn from the $`2^K\left(\genfrac{}{}{0pt}{}{N}{K}\right)`$ possibilities. In the recent years the random $`K`$-satisfiability problem has attracted much interest in computer science and in statistical physics. Its most striking feature is certainly its sharp threshold. Throughout this paper, ‘with high probability’ (w.h.p.) means with a probability which goes to one as $`N\mathrm{}`$. ###### Conjecture 1.1 (Satisfiability Threshold Conjecture) For all $`K2`$, there exists $`\alpha _c(K)`$ such that: * if $`\alpha <\alpha _c(K)`$, $`F_K(N,N\alpha )`$ is satisfiable w.h.p. * if $`\alpha >\alpha _c(K)`$, $`F_K(N,N\alpha )`$ is unsatisfiable w.h.p. The random $`K`$-SAT problem, for $`N`$ large and $`\alpha `$ close to $`\alpha _c(K)`$, provides instances of very hard CNF formulas that can be used as benchmarks for algorithms. For such hard ensembles, the study of the typical complexity could be crucial for the understanding of the usual ‘worst-case’ complexity. Although Conjecture 1.1 remains unproved, Friedgut established the existence of a non-uniform sharp threshold Friedgut . ###### Theorem 1.2 (Friedgut) For each $`K2`$, there exists a sequence $`\alpha _N(K)`$ such that for all $`ϵ>0`$: $$\underset{N\mathrm{}}{lim}𝐏(F_K(N,N\alpha )\text{ is satisfiable})=\{\begin{array}{cc}1\hfill & \text{if }\alpha =(1ϵ)\alpha _N(K)\hfill \\ 0\hfill & \text{if }\alpha =(1+ϵ)\alpha _N(K).\hfill \end{array}$$ (1) A lot of efforts have been devoted to finding tight bounds for the threshold. The best upper bounds so far were derived using first moment methods kirousis ; dubois , and the best lower bounds were obtained by second moment methods achliomoore ; achlioperes . Using these bounds, it was shown that $`\alpha _c(K)=2^K\mathrm{ln}(2)O(K)`$ as $`K\mathrm{}`$. On the other hand, powerful, self-consistent, but non-rigorous tools from statistical physics were used to predict specific values of $`\alpha _c(K)`$, as well as heuristical asymptotic expansions for large $`K`$ MZ ; MPZ ; MMZ-RSA . The cavity method Cavity , which provides these results, relies on several unproven assumptions motivated by spin-glass theory, the most important of which is the partition of the space of SAT-assignments into many states or clusters far away from each other (with Hamming distance greater than $`cN`$ as $`N\mathrm{}`$), in the so-called hard-SAT phase. So far, the existence of such a clustering phase has been shown rigorously in the simpler case of the random XORSAT problem XORSAT-CDMM ; XORSAT-MRZ ; XORSAT-DM in compliance with the prediction of the cavity method, but its existence is predicted in many other problems, such as $`q`$-colorability mulet ; braunstein or the Multi-Index Matching Problem martinmezardrivoire . At the heuristic level, clustering is an important phenomenon, often held responsible for entrapping local search algorithm into non-optimal metastable states montanarisemerjian . It is also a limiting feature for the belief propagation iterative decoding algorithms in Low Density Parity Check Codes montanari ; FLMR . In this paper we provide a rigorous analysis of some geometrical properties of the space of SAT-assignments in the random $`K`$-SAT problem. This study complements the results of MMZ\_prl , and its results are consistent with the clustering scenario. A new characterizing feature of CNF formulas, the ‘$`x`$-satisfiability’, is proposed, which carries information about the spectrum of distances between SAT-assignments. The $`x`$-satisfiability property is studied thoroughly using first and second moment methods previously developed for the satisfiability threshold. The Hamming distance between two assignments $`(\stackrel{}{\sigma },\stackrel{}{\tau })`$ is defined by $$d_{\stackrel{}{\sigma }\stackrel{}{\tau }}=\frac{N}{2}\frac{1}{2}\underset{i=1}{\overset{N}{}}\sigma _i\tau _i.$$ (2) (Throughout the paper the term ‘distance’ will always refer to the Hamming distance.) Given a random formula $`F_K(N,N\alpha )`$, we define a ‘SAT-$`x`$-pair’ as a pair of assignments $`(\stackrel{}{\sigma },\stackrel{}{\tau })\{1,1\}^{2N}`$, which both satisfy $`F`$, and which are at a fixed distance specified by $`x`$ as follows: $$d_{\stackrel{}{\sigma }\stackrel{}{\tau }}[Nxϵ(N),Nx+ϵ(N)].$$ (3) Here $`x`$ is the proportion of distinct values between the two configurations, which we keep fixed as $`N`$ and $`d`$ go to infinity. The resolution $`ϵ(N)`$ has to be $`1`$ and sub-extensive: $`lim_N\mathrm{}ϵ(N)/N=0`$, but its precise form is unimportant for our large $`N`$ analysis. For example we can choose $`ϵ(N)=\sqrt{N}`$. ###### Definition 1.3 A CNF formula is $`x`$-satisfiable if it possesses a SAT-$`x`$-pair. Note that for $`x=0`$, $`x`$-satisfiability is equivalent to satisfiability, while for $`x=1`$, it is equivalent to Not-All-Equal satisfiability, where each clause must contain at least one satisfied literal and at least one unsatisfied litteral gareyjohnson . The clustering property found heuristically in MPZ ; MZ suggests the following: ###### Conjecture 1.4 For all $`KK_0`$, there exist $`\alpha _1(K)`$, $`\alpha _2(K)`$, with $`\alpha _1(K)<\alpha _2(K)`$, such that: for all $`\alpha (\alpha _1(K),\alpha _2(K))`$, there exist $`x_1(K,\alpha )<x_2(K,\alpha )<x_3(K,\alpha )`$ such that: * for all $`x[0,x_1(K,\alpha )][x_2(K,\alpha ),x_3(K,\alpha )]`$, a random formula $`F_K(N,N\alpha )`$ is $`x`$-satisfiable w.h.p. * for all $`x[x_1(K,\alpha ),x_2(K,\alpha )][x_3(K,\alpha ),1]`$, a random formula $`F_K(N,N\alpha )`$ is $`x`$-unsatisfiable w.h.p. Let us give a geometrical interpretation of this conjecture. The space of SAT-assignments is partioned into non-empty regions whose diameter is smaller than $`x_1`$; the distance between any two of these regions is at least $`x_2`$, while $`x_3`$ is the maximum distance between any pair of SAT-assignments. This interpretation is compatible with the notion of clusters used in the statistical physics approach. It should also be mentioned that in a contribution posterior to this work AchlioptasRicci06 , the number of regions was shown to be exponential in the size of the problem, further supporting the statistical mechanics picture. Conjecture 1.4 can be rephrased in a slightly different way, which decomposes it into two steps. The first step is to state the Satisfiability Threshold Conjecture for pairs: ###### Conjecture 1.5 For all $`K2`$ and for all $`x`$, $`0<x<1`$, there exists an $`\alpha _c(K,x)`$ such that: * if $`\alpha <\alpha _c(x)`$, $`F_K(N,N\alpha )`$ is $`x`$-satisfiable w.h.p. * if $`\alpha >\alpha _c(x)`$, $`F_K(N,N\alpha )`$ is $`x`$-unsatisfiable w.h.p. The second step conjectures that for $`K`$ large enough, as a function of $`x`$, the function $`\alpha _c(K,x)`$ is non monotonic and has two maxima: a local maximum at a value $`x_M(K)<1`$, and a global maximum at $`x=0`$. In this paper we prove the equivalent of Friedgut’s theorem: ###### Theorem 1.6 For each $`K3`$ and $`x`$, $`0<x<1`$, there exists a sequence $`\alpha _N(K,x)`$ such that for all $`ϵ>0`$: $$\underset{N\mathrm{}}{lim}𝐏(F_K(N,N\alpha )\text{ is }x\text{-satisfiable})=\{\begin{array}{cc}1\hfill & \text{if }\alpha =(1ϵ)\alpha _N(K,x),\hfill \\ 0\hfill & \text{if }\alpha =(1+ϵ)\alpha _N(K,x),\hfill \end{array}$$ (4) and we obtain two functions, $`\alpha _{LB}(K,x)`$ and $`\alpha _{UB}(K,x)`$, such that: * For $`\alpha >\alpha _{UB}(K,x)`$, a random $`K`$-CNF $`F_K(N,N\alpha )`$ is $`x`$-unsatisfiable w.h.p. * For $`\alpha <\alpha _{LB}(K,x)`$, a random $`K`$-CNF $`F_K(N,N\alpha )`$ is $`x`$-satisfiable w.h.p. The two functions $`\alpha _{LB}(K,x)`$ and $`\alpha _{UB}(K,x)`$ are lower and upper bounds for $`\alpha _N(K,x)`$ as $`N`$ tends to infinity. Numerical computations of these bounds indicate that $`\alpha _N(K,x)`$ is non monotonic as a function of $`x`$ for $`K8`$, as illustrated in Fig. 1. More precisely, we prove ###### Theorem 1.7 For all $`ϵ>0`$, there exists $`K_0`$ such that for all $`KK_0`$, $`\underset{x(0,\frac{1}{2})}{\mathrm{min}}\alpha _{UB}(K,x)`$ $``$ $`(1+ϵ){\displaystyle \frac{2^K\mathrm{ln}2}{2}},`$ (5) $`\alpha _{LB}(K,0)`$ $``$ $`(1ϵ)2^K\mathrm{ln}2,`$ (6) $`\alpha _{LB}(K,1/2)`$ $``$ $`(1ϵ)2^K\mathrm{ln}2.`$ (7) This in turn shows that, for $`K`$ large enough and in some well chosen interval of $`\alpha `$ below the satisfiability threshold $`\alpha _c2^K\mathrm{ln}2`$, SAT-$`x`$-pairs exist for $`x`$ close to zero and for $`x=\frac{1}{2}`$, but they do not exist in the intermediate $`x`$ region. Note that Eq. (6) was established by achlioperes . In section 2 we establish rigorous and explicit upper bounds using the first-moment method. The existence of a gap interval is proven in a certain range of $`\alpha `$, and bounds on this interval are found, which imply Eq. (5) in Theorem 1.7. Section 3 derives the lower bound, using a weighted second-moment method, as developed recently in achliomoore ; achlioperes , and presents numerical results. In section 4 we discuss the behavior of the lower bound for large $`K`$. The case of $`x=\frac{1}{2}`$ is treated rigorously, and Eq. (7) in Theorem 1.7 is proven. Other values of $`x`$ are treated at the heuristic level. Section 5 presents a proof of Theorem 1.6. We discuss our results in section 6. ## 2 Upper bound: the first moment method The first moment method relies on Markov’s inequality: ###### Lemma 2.1 Let $`X`$ be a non-negative random variable. Then $$𝐏(X1)𝐄(X).$$ (8) We take $`X`$ to be the number of pairs of SAT-assignments at fixed distance: $$Z(x,F)=\underset{\stackrel{}{\sigma },\stackrel{}{\tau }}{}\delta \left(d_{\stackrel{}{\sigma }\stackrel{}{\tau }}[Nx+ϵ(N),Nxϵ(N)]\right)\delta \left[\stackrel{}{\sigma },\stackrel{}{\tau }S(F)\right],$$ (9) where $`F=F_K(N,N\alpha )`$ is a random $`K`$-CNF formula, and $`S(F)`$ is the set of SAT-assignments to this formula. Throughout this paper $`\delta (A)`$ is an indicator function, equal to $`1`$ if the statement $`A`$ is true, equal to $`0`$ otherwise. The expectation $`𝐄`$ is over the set of random $`K`$-CNF formulas. Since $`Z(x,F)1`$ is equivalent to ‘$`F`$ is $`x`$-satisfiable’, (8) gives an upper bound for the probability of $`x`$-satisfiability. The expected value of the double sum can be rewritten as: $$𝐄(Z)=2^N\underset{d[Nx+ϵ(N),Nxϵ(N)]}{}\left(\genfrac{}{}{0pt}{}{N}{d}\right)𝐄\left[\delta \left(\stackrel{}{\sigma },\stackrel{}{\tau }S(F)\right)\right].$$ (10) where $`\stackrel{}{\sigma }`$ and $`\stackrel{}{\tau }`$ are any two assignments with Hamming distance $`d`$. We have $`\delta \left(\stackrel{}{\sigma },\stackrel{}{\tau }S(F)\right)=_c\delta \left(\stackrel{}{\sigma },\stackrel{}{\tau }S(c)\right)`$, where $`c`$ denotes one of the $`M`$ clauses. All clauses are drawn independently, so that we have: $$𝐄(Z)(2ϵ(N)+1)2^N\underset{d[Nx+ϵ(N),Nxϵ(N)]}{\mathrm{max}}\left\{\left(\genfrac{}{}{0pt}{}{N}{d}\right)\left(𝐄\left[\delta \left(\stackrel{}{\sigma },\stackrel{}{\tau }S(c)\right)\right]\right)^M\right\},$$ (11) where we have bounded the sum by the maximal term times the number of terms. $`𝐄\left[\delta \left(\stackrel{}{\sigma },\stackrel{}{\tau }S(c)\right)\right]`$ can easily be calculated and its value is: $`12^{1K}+2^K(1x)^K+o(1)`$. Indeed there are only two realizations of the clause among $`2^K`$ that do not satisfy $`c`$ unless the two configurations overlap exactly on the domain of $`c`$. Considering the normalized logarithm of this quantity, $$F(x,\alpha )=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{ln}𝐄(Z)=\mathrm{ln}2+H_2(x)+\alpha \mathrm{ln}\left(12^{1K}+2^K(1x)^K\right),$$ (12) where $`H_2(x)=x\mathrm{ln}x(1x)\mathrm{ln}(1x)`$ is the two-state entropy function, one can deduce an upper bound for $`\alpha _N(K,x)`$. Indeed, $`F(x,\alpha )<0`$ implies $`lim_N\mathrm{}𝐏(Z(x,F)1)=0`$. Therefore: ###### Theorem 2.2 For each $`K`$ and $`0<x<1`$, and for all $`\alpha `$ such that $$\alpha >\alpha _{UB}(K,x)=\frac{\mathrm{ln}2+H_2(x)}{\mathrm{ln}(12^{1K}+2^K(1x)^K)},$$ (13) a random formula $`F_K(N,N\alpha )`$ is $`x`$-unsatisfiable w.h.p. We observe numerically that a ‘gap’ ($`x_1,x_2`$ and $`\alpha `$ such that $`x_1<x<x_2F(x,\alpha )<0`$) appears for $`K6`$. More generally, the following results holds, which implies Eq. (5) in Theorem 1.7: ###### Theorem 2.3 Let $`ϵ(0,1)`$, and $`\{y_K\}_K`$ be a sequence verifying $`Ky_K\mathrm{}`$ and $`y_K=o(1)`$. Denote by $`H_2^1(u)`$ the smallest root to $`H_2(x)=u`$, with $`u[0,\mathrm{ln}2]`$. There exists $`K_0`$ such that for all $`KK_0`$, $`\alpha [(1+ϵ)2^{K1}\mathrm{ln}2,\alpha _N(K))`$ and $`x[y_K,H_2^1(\alpha 2^{1K}\mathrm{ln}2ϵ)][1H_2^1(\alpha 2^{1K}\mathrm{ln}2ϵ),1]`$, $`F_K(N,N\alpha )`$ is $`x`$-unsatisfiable w.h.p. Proof. Clearly $`(1+ϵ)2^{K1}\mathrm{ln}(2)<\alpha _N(K)`$ since $`\alpha _N(K)=2^K\mathrm{ln}(2)O_K(K)`$ achlioperes . Observe that $`(1y_K)^K=o(1)`$. Then for all $`\delta >0`$, there exists $`K_1`$ such that for all $`KK_1`$, $`x>y_K`$: $$\alpha _{UB}(x)<(1+\delta )2^{K1}(\mathrm{ln}2+H_2(x)).$$ (14) Inverting this inequality yields the theorem. $`\mathrm{}`$ The choice (9) of $`X`$, although it is the simplest one, is not optimal. The first moment method only requires the condition $`X1`$ to be equivalent to the $`x`$-satisfiability, and better choices of $`X`$ exist which allow to improve the bound. Techniques similar to the one introduced separately by Dubois and Boufkhad dubois on the one hand, and Kirousis, Kranakis and Krizanc kirousis on the other hand, can be used to obtain two tighter bounds. Quantitatively, it turns out that these more elaborate bounds provide only very little improvement over the simple bound (13) (see Fig. 2). For the sake of completeness, we give without proof the simplest of these bounds: ###### Theorem 2.4 The unique positive solution of the equation $`H_2(x)`$ $`+\alpha \mathrm{ln}\left(12^{1K}+2^K(1x)^K\right)`$ $`+(1x)\mathrm{ln}\left[2\mathrm{exp}\left(K\alpha {\displaystyle \frac{2^{1K}2^K(1x)^{K1}}{12^{1K}+2^K(1x)^K}}\right)\right]`$ $`+x\mathrm{ln}\left[2\mathrm{exp}\left(K\alpha {\displaystyle \frac{2^{1K}2^{1K}(1x)^{K1}}{12^{1K}+2^K(1x)^K}}\right)\right]=0`$ (15) is an upper bound for $`\alpha _N(K,x)`$. For $`x=0`$ we recover the expression of kirousis . The proof closely follows that of kirousis and presents no notable difficulty. We also derived a tighter bound based on the technique used in dubois , gaining only a small improvement over the bound of Theorem 2.4 (less than $`.001\%`$). ## 3 Lower bound: the second moment method The second moment method uses the following consequence of Chebyshev’s inequality: ###### Lemma 3.1 If $`X`$ is a non-negative random variable, one has: $$𝐏(X>0)\frac{𝐄(X)^2}{𝐄(X^2)}.$$ (16) It is well known that the simplest choice of $`X`$ as the number of SAT-assignments (in our case the number of SAT-$`x`$-pairs) is bound to fail. The intuitive reason achliomoore ; achlioperes is that this naive choice favors pairs of SAT-assignments with a great number of satisfying litterals. It turns out that such assignments are highly correlated, since they tend to agree with each other, and this causes the failure of the second-moment method. In order to deal with balanced (with approximately half of literals satisfied) and uncorrelated pairs of assignments, one must consider a weighted sum of all SAT-assignments. Following achliomoore ; achlioperes , we define: $$Z(x,F)=\underset{\stackrel{}{\sigma },\stackrel{}{\tau }}{}\delta \left(d_{\stackrel{}{\sigma }\stackrel{}{\tau }}=Nx\right)W(\stackrel{}{\sigma },\stackrel{}{\tau },F),$$ (17) where $`Nx`$ denotes the integer part of $`Nx`$. Note that the condition $`d_{\stackrel{}{\sigma }\stackrel{}{\tau }}=Nx`$ is stronger than Eq. (3). The weights $`W(\stackrel{}{\sigma },\stackrel{}{\tau },F)`$ are decomposed according to each clause: $`W(\stackrel{}{\sigma },\stackrel{}{\tau },F)`$ $`=`$ $`{\displaystyle \underset{c}{}}W(\stackrel{}{\sigma },\stackrel{}{\tau },c),`$ (18) $`\text{with}W(\stackrel{}{\sigma },\stackrel{}{\tau },c)`$ $`=`$ $`W(\stackrel{}{u},\stackrel{}{v}),`$ (19) where $`\stackrel{}{u},\stackrel{}{v}`$ are $`K`$-component vectors such that: $`u_i=1`$ if the $`i^{\text{th}}`$ litteral of $`c`$ is satisfied under $`\stackrel{}{\sigma }`$, and $`u_i=1`$ otherwise (here we assume that the variables connected to $`c`$ are arbitrarily ordered). $`\stackrel{}{v}`$ is defined in the same way with respect to $`\stackrel{}{\tau }`$. In order to have the equivalence between $`Z>0`$ and the existence of pairs of SAT-assignments, we impose the following condition on the weights: $$W(\stackrel{}{u},\stackrel{}{v})=\{\begin{array}{cc}0\hfill & \text{if}\stackrel{}{u}=(1,\mathrm{},1)\text{or}\stackrel{}{v}=(1,\mathrm{},1),\hfill \\ >0\hfill & \text{otherwise}.\hfill \end{array}$$ (20) Let us now compute the first and second moments of $`Z`$: ###### Claim 3.2 $$𝐄(Z)=2^N\left(\genfrac{}{}{0pt}{}{N}{Nx}\right)f_1(x)^M,$$ (21) where $`f_1(x)`$ $`=`$ $`𝐄[W(\stackrel{}{\sigma },\stackrel{}{\tau },c)]`$ (22) $`=`$ $`2^K{\displaystyle \underset{\stackrel{}{u},\stackrel{}{v}}{}}W(\stackrel{}{u},\stackrel{}{v})(1x)^{\|\stackrel{}{uv}\|}x^{K\|\stackrel{}{uv}\|}.`$ (23) Here $`\|\stackrel{}{u}\|`$ is the number of indices $`i`$ such that $`u_i=+1`$, and $`\stackrel{}{uv}`$ denotes the vector $`(u_1v_1,\mathrm{},u_Kv_K)`$. Writing the second moment is a little more cumbersome: ###### Claim 3.3 $$𝐄(Z^2)=2^N\underset{𝐚V_N\{0,1/N,2/N,\mathrm{},1\}^8}{}\frac{N!}{_{i=0}^7(Na_i)!}f_2(𝐚)^M,$$ (24) where $`f_2(𝐚)`$ $`=`$ $`𝐄[W(\stackrel{}{\sigma },\stackrel{}{\tau },c)W(\stackrel{}{\sigma },\stackrel{}{\tau },c)]`$ (25) $`=`$ $`2^K{\displaystyle \underset{\stackrel{}{u},\stackrel{}{v},\stackrel{}{u}^{},\stackrel{}{v}^{}}{}}W(\stackrel{}{u},\stackrel{}{v})W(\stackrel{}{u}^{},\stackrel{}{v}^{}){\displaystyle \underset{i=1}{\overset{K}{}}}a_0^{\delta (u_i=v_i=u_i^{}=v_i^{})}a_1^{\delta (u_i=v_i=u_i^{}v_i^{})}`$ $`a_2^{\delta (u_i=v_i=v_i^{}u_i^{})}a_3^{\delta ((u_i=v_i)(u_i^{}=v_i^{}))}a_4^{\delta (u_i=u_i^{}=v_i^{}v_i)}`$ $`a_5^{\delta ((u_i=u_i^{})(v_i=v_i^{}))}a_6^{\delta ((u_i=v_i^{})(u_i^{}=v_i))}a_7^{\delta (u_i^{}=v_i^{}=u_iu_i)}`$ $`𝐚`$ is a 8-component vector giving the proportion of each type of quadruplets $`(\tau _i,\sigma _i,\tau _i^{},\sigma _i^{})`$ — $`\stackrel{}{\tau }`$ being arbitrarily (but without losing generality) fixed to $`(1,\mathrm{},1)`$ — as described in the following table: \| \| $`a_0`$ \| $`a_1`$ \| $`a_2`$ \| $`a_3`$ \| $`a_4`$ \| $`a_5`$ \| $`a_6`$ \| $`a_7`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`\tau _i`$ \| + \| + \| + \| + \| + \| + \| + \| + \| \| $`\sigma _i`$ \| + \| + \| + \| + \| $``$ \| $``$ \| $``$ \| $``$ \| \| $`\tau _i^{}`$ \| + \| + \| $``$ \| $``$ \| + \| + \| $``$ \| $``$ \| \| $`\sigma _i^{}`$ \| + \| $``$ \| + \| $``$ \| + \| $``$ \| + \| $``$ \| The set $`V_N[0,1]^8`$ is a simplex specified by: $$\{\begin{array}{c}N(a_4+a_5+a_6+a_7)=Nx\hfill \\ N(a_1+a_2+a_5+a_6)=Nx\hfill \\ _{i=0}^7a_i=1\hfill \end{array}$$ (26) These three conditions (26) correspond to the normalization of the proportions and to the enforcement of the conditions $`d_{\stackrel{}{\sigma }\stackrel{}{\tau }}=Nx`$, $`d_{\stackrel{}{\sigma }^{}\stackrel{}{\tau }^{}}=Nx`$. When $`N\mathrm{}`$, $`V=_NV_N`$ defines a five-dimensional simplex described by the three hyperplanes: $$\{\begin{array}{c}a_4+a_5+a_6+a_7=x\hfill \\ a_1+a_2+a_5+a_6=x\hfill \\ _{i=0}^7a_i=1\hfill \end{array}$$ (27) In order to yield an asymptotic estimate of $`𝐄(Z^2)`$ we first use the following lemma, which results from a simple approximation of integrals by sums: ###### Lemma 3.4 Let $`\psi (𝐚)`$ be a real, positive, continuous function of $`𝐚`$, and let $`V_N`$, $`V`$ be defined as previously. Then there exists a constant $`C_0`$ depending on $`x`$ such that for sufficiently large $`N`$: $$\underset{𝐚V_N\{1/N,2/N,\mathrm{},1\}^8}{}\frac{N!}{_{i=0}^7(Na_i)!}\psi (𝐚)^NC_0N^{3/2}_Vd𝐚e^{N[H_8(𝐚)+\mathrm{ln}\psi (𝐚)]},$$ (28) where $`H_8(𝐚)=_{i=1}^8a_i\mathrm{ln}a_i`$. A standard Laplace method used on Eq. (28) with $`\psi =2(f_2)^\alpha `$ yields: ###### Claim 3.5 For each $`K,x`$, define: $$\mathrm{\Phi }(𝐚)=H_8(𝐚)\mathrm{ln}22H_2(x)+\alpha \mathrm{ln}f_2(𝐚)2\alpha \mathrm{ln}f_1(x).$$ (29) and let $`𝐚_0V`$ be the global maximum of $`\mathrm{\Phi }`$ restricted to $`V`$. Suppose that $`_𝐚^2\mathrm{\Phi }(𝐚_0)`$ is definite negative. Then there exists a constant $`C_1`$ such that, for $`N`$ sufficiently large, $$\frac{𝐄(Z)^2}{𝐄(Z^2)}C_1\mathrm{exp}(N\mathrm{\Phi }(𝐚_0)).$$ (30) Obviously $`\mathrm{\Phi }(𝐚_0)0`$ in general. In order to use Lemma 3.1, one must find the weights $`W(\stackrel{}{u},\stackrel{}{v})`$ in such a way that $`\mathrm{max}_{𝐚V}\mathrm{\Phi }(𝐚)=0`$. We first notice that, at the particular point $`𝐚^{}`$ where the two pairs are uncorrelated with each other, $$a_0^{}=a_3^{}=\frac{(1x)^2}{2},a_1^{}=a_2^{}=a_4^{}=a_7^{}=\frac{x(1x)}{2},a_5^{}=a_6^{}=\frac{x^2}{2},$$ (31) we have the following properties: * $`H_8(𝐚^{})=\mathrm{ln}2+2H_2(x)`$, * $`_𝐚H_8(𝐚^{})=0,`$ $`_𝐚^2H_8(𝐚^{})`$ definite negative, * $`f_1(x)^2=f_2(𝐚^{})`$ and hence $`\mathrm{\Phi }(𝐚^{})=0`$. (Note that the derivatives $`_𝐚`$ are taken in the simplex $`V`$). So the weights must be chosen in such a way that $`𝐚^{}`$ be the global maximum of $`\mathrm{\Phi }`$. A necessary condition is that $`𝐚^{}`$ be a local maximum, which entails $`_𝐚f_2(𝐚^{})=0`$. Using the fact that the number of common values between four vectors $`\stackrel{}{u},\stackrel{}{v},\stackrel{}{u}^{},\stackrel{}{v}^{}\{1,1\}^K`$ can be written as: $$\frac{1}{8}\left(K+\stackrel{}{u}\stackrel{}{v}+\stackrel{}{u}\stackrel{}{u}^{}+\stackrel{}{u}\stackrel{}{v}^{}+\stackrel{}{v}\stackrel{}{u}^{}+\stackrel{}{v}\stackrel{}{v}^{}+\stackrel{}{u}^{}\stackrel{}{v}^{}+\stackrel{}{uv}\stackrel{}{u^{}v^{}}\right)$$ (32) we deduce from $`_𝐚f_2(𝐚^{})=0`$ the condition: $$\underset{\stackrel{}{u},\stackrel{}{v}}{}W(\stackrel{}{u},\stackrel{}{v})\{\begin{array}{c}\stackrel{}{u}\hfill \\ \stackrel{}{v}\hfill \end{array}(1x)^{\|\stackrel{}{uv}\|}x^{K\|\stackrel{}{uv}\|}=0,$$ (33) $`0`$ $`=`$ $`K(2x1)^2\left[{\displaystyle \underset{\stackrel{}{u},\stackrel{}{v}}{}}W(\stackrel{}{u},\stackrel{}{v})(1x)^{\|\stackrel{}{uv}\|}x^{K\|\stackrel{}{uv}\|}\right]^2`$ (34) $`+\left[{\displaystyle \underset{\stackrel{}{u},\stackrel{}{v}}{}}W(\stackrel{}{u},\stackrel{}{v})\stackrel{}{uv}(1x)^{\|\stackrel{}{uv}\|}x^{K\|\stackrel{}{uv}\|}\right]^2`$ $`+2(2x1)\left[{\displaystyle \underset{\stackrel{}{u},\stackrel{}{v}}{}}W(\stackrel{}{u},\stackrel{}{v})\stackrel{}{u}\stackrel{}{v}(1x)^{\|\stackrel{}{uv}\|}x^{K\|\stackrel{}{uv}\|}\right]`$ $`\times \left[{\displaystyle \underset{\stackrel{}{u},\stackrel{}{v}}{}}W(\stackrel{}{u},\stackrel{}{v})(1x)^{\|\stackrel{}{uv}\|}x^{K\|\stackrel{}{uv}\|}\right].`$ If we suppose that $`W`$ is invariant under simultaneous and identical permutations of the $`u_i`$ or of the $`v_i`$ (which we must, since the ordering of the variables by the label $`i`$ is arbitrary), the $`K`$ components of all vectorial quantities in Eqs. (33), (34) should be equal. Then we obtain equivalently: $`{\displaystyle \underset{\stackrel{}{u},\stackrel{}{v}}{}}W(\stackrel{}{u},\stackrel{}{v})(2\|\stackrel{}{u}\|K)(1x)^{\|\stackrel{}{uv}\|}x^{K\|\stackrel{}{uv}\|}=0\text{and}\stackrel{}{u}\stackrel{}{v},`$ (35) $`{\displaystyle \underset{\stackrel{}{u},\stackrel{}{v}}{}}W(\stackrel{}{u},\stackrel{}{v})(K(2x1)+\stackrel{}{u}\stackrel{}{v})(1x)^{\|\stackrel{}{uv}\|}x^{K\|\stackrel{}{uv}\|}=0,`$ (36) We choose the following simple form for $`W(\stackrel{}{u},\stackrel{}{v})`$: $$W(\stackrel{}{u},\stackrel{}{v})=\{\begin{array}{cc}0\hfill & \text{if}\stackrel{}{u}=(1,\mathrm{},1)\text{or}\stackrel{}{v}=(1,\mathrm{},1),\hfill \\ \lambda ^{\|\stackrel{}{u}\|+\|\stackrel{}{v}\|}\nu ^{\|\stackrel{}{uv}\|}\hfill & \text{otherwise}.\hfill \end{array}$$ (37) Although this choice is certainly not optimal, it turns out particularly tractable. Eqs. (35) and (36) simplify to: $$\begin{array}{cc}\hfill [\nu (1x)]^{K1}=& (\lambda ^2+12\lambda \nu )\left(2\lambda x+\nu (1x)(1+\lambda ^2)\right)^{K1}\hfill \\ \hfill \left(\nu (1x)+\lambda x\right)^{K1}=& (1\lambda \nu )\left(2\lambda x+\nu (1x)(1+\lambda ^2)\right)^{K1}.\hfill \end{array}$$ (38) We found numerically a unique solution $`\lambda >0,\nu >0`$ to these equations for any value of $`K2`$ that we checked. Fixing $`(\lambda ,\nu )`$ to a solution of (38), we seek the largest value of $`\alpha `$ such that the local maximum $`𝐚^{}`$ is a global maximum, i.e. such that there exists no $`𝐚V`$ with $`\mathrm{\Phi }(𝐚)>0`$. To proceed one needs analytical expressions for $`f_1(x)`$ and $`f_2(𝐚)`$. $`f_1`$ simply reads: $`f_1(x)`$ $`=`$ $`2^K\left((1x)\nu (1+\lambda ^2)+2x\lambda \right)^K22^K\left(x\lambda +(1x)\nu \right)^K`$ (39) $`+2^K((1x)\nu )^K.`$ $`f_2`$ is calculated by Sylvester’s formula, but its expression is long and requires preliminar notations. We index the $`16`$ possibilities for $`(u_i,v_i,u_i^{},v_i^{})`$ by a number $`r\{0,\mathrm{},15\}`$ defined as: $$r=8\frac{1u_i}{2}+4\frac{1v_i}{2}+2\frac{1u_i^{}}{2}+\frac{1v_i^{}}{2}.$$ (40) For each index $`r`$, define $`l(r)`$ $`=\delta (u_i=1)+\delta (v_i=1)+\delta (u_i^{}=1)+\delta (v_i^{}=1),`$ (41) $`n(r)`$ $`=\delta (u_iv_i=1)+\delta (u_i^{}v_i^{}=1),`$ (42) and $`z_r`$ $`=`$ $`\lambda ^{l(r)}\nu ^{n(r)}\times \{\begin{array}{cc}a_r\hfill & \text{if}r7\hfill \\ a_{15r}\hfill & \text{if}r8\hfill \end{array}.`$ (45) Also define the four following subsets of $`\{0,\mathrm{},15\}`$: $`A_0`$ is the set of indices $`r`$ corresponding to quadruplets of the form $`(1,v_i,u_i^{},v_i^{})`$. $`A_0=\{r\{0,\mathrm{},15\}\|u_i=1\}`$. Similarly, $`A_1=\{r\|v_i=1\}`$, $`A_2=\{r\|u_i^{}=1\}`$ and $`A_3=\{r\|v_i^{}=1\}`$. Then $`f_2`$ is given by: $`2^Kf_2(𝐚)`$ $`=`$ $`\left({\displaystyle \underset{j=0}{\overset{15}{}}}z_j\right)^K{\displaystyle \underset{k=0}{\overset{3}{}}}\left({\displaystyle \underset{jA_k}{}}z_j\right)^K+{\displaystyle \underset{0k<k^{}3}{}}\left({\displaystyle \underset{jA_kA_k^{}}{}}z_j\right)^K`$ (46) $`{\displaystyle \underset{0k<k^{}<k^{\prime \prime }3}{}}\left({\displaystyle \underset{jA_kA_k^{}A_{k^{\prime \prime }}}{}}z_j\right)^K+\left({\displaystyle \underset{jA_0A_1A_2A_3}{}}z_j\right)^K.`$ We can now state our lower-bound result: ###### Lemma 3.6 Let $`\alpha _+(0,+\mathrm{}]`$ be the smallest $`\alpha `$ such that $`_𝐚^2\mathrm{\Phi }(𝐚^{})`$ is not definite negative. For each $`K`$ and $`x(0,1)`$, and for all $`\alpha \alpha _{LB}(K,x)`$, with $$\alpha _{LB}(K,x)=\mathrm{min}[\alpha _+,\underset{𝐚V_+}{inf}\frac{\mathrm{ln}2+2H_2(x)H_8(𝐚)}{\mathrm{ln}f_2(𝐚)2\mathrm{ln}f_1(x)}],$$ (47) where $`V_+=\{𝐚V\|f_2(𝐚)>f_1^2\left(1/2\right)\}`$, and where $`(\lambda ,\nu )`$ is chosen to be a positive solution of (38), the probability that a random formula $`F_K(N,N\alpha )`$ is $`x`$-satisfiable is bounded away from $`0`$ as $`N\mathrm{}`$. This is a straightforward consequence of the expression (29) of $`\mathrm{\Phi }(𝐚)`$. Theorem 1.6 and Lemma 3.6 immediately imply: ###### Theorem 3.7 For all $`\alpha <\alpha _{LB}(K,x)`$ defined in Lemma 3.6, a random $`K`$-CNF formula $`F_K(N,N\alpha )`$ is $`x`$-satisfiable w.h.p. We devised several numerical strategies to evaluate $`\alpha _{LB}(K,x)`$. The implementation of Powell’s method on each point of a grid of size $`𝒩^5`$ ($`𝒩=10,15,20`$) on $`V`$ turned out to be the most efficient and reliable. The results are given by Fig. 1 for $`K=8`$, the smallest $`K`$ such that the picture given by Conjecture 1.4 is confirmed. We found a clustering phenomenon for all the values of $`K8`$ that we checked. In the following we shall provide a rigorous estimate of $`\alpha _{LB}(K,\frac{1}{2})`$ at large $`K`$. ## 4 Large $`K`$ analysis ### 4.1 Asymptotics for $`x=\frac{1}{2}`$ The main result of this section is contained in the following theorem, which implies Eq. (7) in Theorem 1.7: ###### Theorem 4.1 The large $`K`$ asymptotics of $`\alpha _{LB}(K,x)`$ at $`x=1/2`$ is given by: $$\alpha _{LB}(K,1/2)2^K\mathrm{ln}2.$$ (48) The proof primarily relies on the following results: ###### Claim 4.2 Let $`\nu =1`$ and $`\lambda `$ be the unique positive root of: $$(1\lambda )(1+\lambda )^{K1}1=0.$$ (49) Then $`(\lambda ,\nu )`$ is solution to (38) with $`x=\frac{1}{2}`$ and one has, at large $`K`$: $$\lambda 12^{1K}.$$ (50) ###### Lemma 4.3 Let $`x=\frac{1}{2}`$. There exist $`K_0>0`$, $`C_1>0`$ and $`C_2>0`$ such that for all $`KK_0`$, and for all $`𝐚V`$ s.t. $`\|𝐚𝐚^{}\|<1/8`$, $$\left\|\mathrm{ln}f_2(𝐚)2\mathrm{ln}f_1(1/2)\right\|K^2C_1\|𝐚𝐚^{}\|^22^{2K}+C_2\|𝐚𝐚^{}\|^32^K$$ (51) ###### Lemma 4.4 Let $`x=\frac{1}{2}`$. There exist $`K_0>0`$, $`C_0>0`$ such that for $`KK_0`$, for all $`𝐚V`$, $$\begin{array}{c}\hfill \|\mathrm{ln}f_2(𝐚)2\mathrm{ln}f_1(1/2)\|2^K[(a_0+a_1+a_4+a_5)^K+(a_0+a_2+a_4+a_6)^K\\ \hfill +(a_0+a_1+a_6+a_7)^K+(a_0+a_2+a_5+a_7)^K]+C_0K2^{2K}\end{array}$$ (52) The proofs of these lemmas are defered to sections 4.3 and 4.4. ### 4.2 Proof of Theorem 4.1 We first show that $`_𝐚^2\mathrm{\Phi }(𝐚^{})`$ is definite negative for all $`\alpha <2^K`$, when $`K`$ is sufficiently large. Indeed $`_𝐚^2H_8(𝐚^{})`$ is definite negative and its largest eigenvalue is $`4`$. Using Lemma 4.3, for $`𝐚V`$ close enough to $`𝐚^{}`$: $$\mathrm{\Phi }(𝐚)2\|𝐚𝐚^{}\|^2+\alpha C_1\|𝐚𝐚^{}\|^2K^22^{2K}+\alpha C_2\|𝐚𝐚^{}\|^32^K.$$ (53) Therefore $$\mathrm{\Phi }(𝐚)\|𝐚𝐚^{}\|^2\text{ for }K\text{ large enough, }\|𝐚𝐚^{}\|<\frac{1}{2C_2}\text{ and }\alpha <2^K.$$ (54) Using Theorem 3.6, we need to find the minimum, for $`aV_+`$, of $$G(K,𝐚)\frac{3\mathrm{ln}2H_8(𝐚)}{\mathrm{ln}f_2(𝐚)2\mathrm{ln}f_1(1/2)}.$$ (55) We shall show that $$\underset{𝐚V_+}{inf}G(K,𝐚)2^K\mathrm{ln}2.$$ (56) We divide this task in two parts. The first part states that there exists $`R>0`$ and $`K_1`$ such that for all $`KK_1`$, and for all $`𝐚V_+`$ such that $`\|𝐚𝐚^{}\|<R`$, $`G(K,𝐚)>2^K`$. This is a consequence of Lemma 4.3; using the fact that $`3\mathrm{ln}2H_8(𝐚)\|𝐚𝐚^{}\|^2`$ for $`𝐚`$ close enough to $`𝐚^{}`$, one obtains: $$G(K,𝐚)\frac{2^K}{C_1K^22^K+C_2\|𝐚𝐚^{}\|}$$ (57) which, for $`K`$ large enough and $`𝐚`$ close enough to $`𝐚^{}`$, is greater than $`2^K`$. The second part deals with the case where $`𝐚`$ is far from $`𝐚^{}`$, i.e. $`\|𝐚𝐚^{}\|>R`$. First we put a bound on the numerator of $`G(𝐚)`$: there exists a constant $`C_3>0`$ such that for all $`𝐚V`$ s.t. $`\|𝐚𝐚^{}\|>R`$, one has $`3\mathrm{ln}2H_8(𝐚)>C_3`$. Looking at Eq. (52), it is clear that, in order to minimize $`G(K,𝐚)`$, $`𝐚`$ should be ‘close’ to at least one the four hyperplanes defined by $$\begin{array}{cc}\hfill a_0+a_1+a_4+a_5=1,& a_0+a_2+a_4+a_6=1,\hfill \\ \hfill a_0+a_1+a_6+a_7=1,& a_0+a_2+a_5+a_7=1.\hfill \end{array}$$ (58) More precisely, we say for instance that $`𝐚`$ is close to the first hyperplane defined above iff $$a_0+a_1+a_4+a_5>1K^{1/2}$$ (59) Now suppose that $`𝐚`$ is not close to that hyperplane. Then the corresponding term goes to $`0`$: $$(a_0+a_1+a_4+a_5)^K\left(1K^{1/2}\right)^K\mathrm{exp}(\sqrt{K})\text{as }K\mathrm{}.$$ (60) We classify all possible cases according to the number of hyperplanes $`𝐚V_+`$ is close to: * $`𝐚`$ is close to none of the hyperplanes. Then $$G(K,𝐚)\frac{2^KC_3}{4\mathrm{exp}(\sqrt{K})+C_0K2^K}>2^K\text{for }K\text{ large enough.}$$ (61) * $`𝐚`$ is close to one hyperplane only, e.g. the first hyperplane $`a_0+a_1+a_4+a_5=1`$ (the other hyperplanes are treated equivalently). As $`_{i=0}^7a_i=0`$, one has $$a_2<K^{1/2},a_3<K^{1/2},a_6<K^{1/2},a_7<K^{1/2}.$$ (62) This implies $`H_8(𝐚)<2\mathrm{ln}2+2\mathrm{ln}K/\sqrt{K}`$, and we get: $$G(K,𝐚)\frac{2^K[\mathrm{ln}22\mathrm{ln}K/\sqrt{K}]}{1+C_0K2^K+3\mathrm{}^\sqrt{K}}2^K(\mathrm{ln}2)\left[13\mathrm{ln}K/\sqrt{K}\right]$$ (63) for sufficiently large $`K`$. * $`𝐚`$ is close to two hyperplanes. It is easy to check that these hyperplanes must be either the first and the fourth ones, or the second and the third ones. In the first case we have $`a_0+a_5>13/\sqrt{K}`$ and in the second case $`a_0+a_6>13/\sqrt{K}`$. Both cases imply: $`H_8(𝐚)<\mathrm{ln}2+3\mathrm{ln}K/\sqrt{K}`$. One thus obtains: $$G(K,𝐚)\frac{2^K[2\mathrm{ln}23\mathrm{ln}K/\sqrt{K}]}{2+C_0K2^K+2\mathrm{}^\sqrt{K}}2^K(\mathrm{ln}2)\left[13\mathrm{ln}K/\sqrt{K}\right].$$ (64) * One can check that $`𝐚`$ cannot be close to more than two hyperplanes. To sum up, we have proved that for $`K`$ large enough, for all $`𝐚V_+`$, $$G(K,𝐚)2^K(\mathrm{ln}2)\left[13\mathrm{ln}K/\sqrt{K}\right],$$ (65) Clearly, $`\alpha _{LB}(K,1/2)=inf_{𝐚V_+}G(K,𝐚)<\alpha _{UB}(K,1/2)`$. Since from Theorem 2.2 we know that $`\alpha _{UB}(K,1/2)2^K\mathrm{ln}2`$, this proves Eq. (56). ### 4.3 Proof of Lemma 4.3 Let $`x=\frac{1}{2}`$ and choose $`\nu =1`$ and $`\lambda `$ the unique positive root of Eq. (49). Let $`ϵ_i=a_i1/8`$, and $`\mathit{ϵ}=(ϵ_0,\mathrm{},ϵ_7)`$. We expand $`f_2(𝐚)`$ in series of $`\mathit{ϵ}`$. The zeroth order term is $`f_2(1/8,\mathrm{},1/8)=f_1^2(1/2)`$. The first order term vanishes. We thus get: $$f_2(𝐚)=f_1^2(1/2)+B_0B_1+B_2B_3+B_4,$$ (66) with $`B_0`$ $`=`$ $`{\displaystyle \underset{q=2}{\overset{K}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{K}{q}}\right)\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{i=0}{\overset{7}{}}}p_i(\lambda )ϵ_i\right)^q\left[{\displaystyle \frac{1+\lambda }{2}}\right]^{4(Kq)},`$ (67) $`B_1`$ $`=`$ $`2^K{\displaystyle \underset{a=1}{\overset{4}{}}}{\displaystyle \underset{q=2}{\overset{K}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{K}{q}}\right)\left[{\displaystyle \underset{i=0}{\overset{7}{}}}\left(\lambda ^{\mathrm{}_{ai}}1\right)ϵ_i\right]^q\left[{\displaystyle \frac{1+\lambda }{2}}\right]^{3(Kq)},`$ (68) $`B_2`$ $`=`$ $`2^{2K}{\displaystyle \underset{a=1}{\overset{6}{}}}{\displaystyle \underset{q=2}{\overset{K}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{K}{q}}\right)[2r_a(\lambda ,\mathit{ϵ})]^q\left[{\displaystyle \frac{1+\lambda }{2}}\right]^{2(Kq)},`$ (69) $`B_3`$ $`=`$ $`2^{3K}{\displaystyle \underset{a=1}{\overset{4}{}}}{\displaystyle \underset{q=2}{\overset{K}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{K}{q}}\right)[4s_a(\lambda ,\mathit{ϵ})]^q\left[{\displaystyle \frac{1+\lambda }{2}}\right]^{Kq},`$ (70) $`B_4`$ $`=`$ $`2^{4K}{\displaystyle \underset{k=2}{\overset{K}{}}}(8ϵ_0)^q.`$ (71) In $`B_0`$, $`p_i(\lambda )=\lambda ^{l(i)}+\lambda ^{l(15i)}24(\lambda 1)`$. We have used the fact that $`_{i=0}^7ϵ_i=0`$. Using $`l(i)+l(15i)=4`$, one obtains $`\|p_i(\lambda )\|11(\lambda 1)^2112^{42K}`$, since $`\|\lambda 1\|2^{2K}`$ for $`K`$ large enough, by virtue of Lemma 4.2. In $`B_1`$, we have used again $`_{i=0}^7ϵ_i=0`$. $`\mathrm{}_{ai}`$ is either $`l(i)`$ or $`l(15i)`$, depending on $`a`$. In both cases $`\|\lambda ^{\mathrm{}_{ai}}1\|4\|\lambda 1\|2^{4K}`$. In $`B_2`$ and $`B_3`$, the expressions of $`r_a(\lambda ,\mathit{ϵ})`$ and $`s_a(\lambda ,\mathit{ϵ})`$ are given by: $$\begin{array}{cc}\hfill r_1=ϵ_0+\lambda (ϵ_1+ϵ_2)+\lambda ^2ϵ_3,& r_2=ϵ_0+\lambda (ϵ_1+ϵ_4)+\lambda ^2ϵ_5,\hfill \\ \hfill r_3=ϵ_0+\lambda (ϵ_2+ϵ_4)+\lambda ^2ϵ_6,& r_4=ϵ_0+\lambda (ϵ_1+ϵ_7)+\lambda ^2ϵ_6,\hfill \\ \hfill r_5=ϵ_0+\lambda (ϵ_2+ϵ_7)+\lambda ^2ϵ_5,& r_6=ϵ_0+\lambda (ϵ_4+ϵ_7)+\lambda ^2ϵ_3,\hfill \end{array}$$ (72) $$s_1=ϵ_0+\lambda ϵ_1,s_2=ϵ_0+\lambda ϵ_2,s_3=ϵ_0+\lambda ϵ_4,s_4=ϵ_0+\lambda ϵ_7.$$ (73) In order to prove Lemma 4.3 we will use the following fact: ###### Claim 4.5 Let $`y`$ be a real variable such that $`\|y\|1`$. Then $$\left\|\underset{k=2}{\overset{K}{}}\left(\genfrac{}{}{0pt}{}{K}{k}\right)y^k\right\|\frac{K(K1)}{2}y^2+2^K\|y\|^3.$$ (74) One has $`\|2r_a\|8\|\mathit{ϵ}\|`$, $`\|4s_a\|8\|\mathit{ϵ}\|`$, and $`\|8ϵ_0\|8\|\mathit{ϵ}\|`$. Therefore, for $`\|\mathit{ϵ}\|<1/8`$, one can write: $`\|B_0\|`$ $``$ $`{\displaystyle \frac{K(K1)}{2}}(112^6)^22^{4K}\|\mathit{ϵ}\|^2+(112^6)^32^{5K}\|\mathit{ϵ}\|^3`$ (75) $`\|B_1\|`$ $``$ $`4{\displaystyle \frac{K(K1)}{2}}2^{14}2^{3K}\|\mathit{ϵ}\|^2+2^{21}2^{3K}\|\mathit{ϵ}\|^3`$ (76) $`\|B_i\|`$ $``$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{4}{i}}\right){\displaystyle \frac{K(K1)}{2}}2^62^{iK}\|\mathit{ϵ}\|^2+2^92^{(i1)K}\|\mathit{ϵ}\|^3\text{for }2i4.`$ (77) Observe that $$f_1(1/2)=\left[\left(\frac{1+\lambda }{2}\right)^K2^K\right]^2=1+O(K2^K)$$ (78) and that for $`K`$ large enough, $$\left\|\mathrm{ln}\frac{f_2(𝐚)}{f_1^2(1/2)}\right\|\frac{2}{f_1(1/2)^2}\underset{i=0}{\overset{4}{}}\|B_i\|,$$ (79) which proves Lemma 4.3. ### 4.4 Proof of Lemma 4.4 Note that the bounds on $`B_0`$ and $`B_1`$ (75), (76) remain valid for any $`\mathit{ϵ}`$. Therefore $`B_0=O(2^{2K})`$ and $`B_1=O(2^{2K})`$ uniformly. We bound $`B_3`$ by observing that: $$\begin{array}{cc}\hfill B_3=& 2^K\left[(a_0+\lambda a_1)^K+(a_0+\lambda a_2)^K+(a_0+\lambda a_4)^K+(a_0+\lambda a_7)^K\right]\hfill \\ & 2^{3K}\underset{a=1}{\overset{4}{}}\left[\frac{1+\lambda }{2}\right]^K\left[1+K\left(\frac{8s_a(\lambda ,\mathit{ϵ})}{1+\lambda }\right)\right].\hfill \end{array}$$ (80) Since $`(a_0+\lambda a_1)a_0+a_11/2`$ and likewise for the three other terms, one has $`B_3=O(2^{2K})`$ uniformly in $`𝐚`$. A similar argument yields $`B_4=O(2^{2K})`$. There remains $`B_2`$, which we write as: $$\begin{array}{cc}\hfill B_2=& 2^K\underset{0k<k^{}3}{}\left(\underset{jA_kA_k^{}}{}z_j\right)^K\hfill \\ & 2^{2K}\underset{a=1}{\overset{6}{}}\left[\frac{1+\lambda }{2}\right]^{2K}\left[1+K\left(\frac{8r_a(\lambda ,\mathit{ϵ})}{(1+\lambda )^2}\right)\right]\hfill \end{array}$$ (81) The second term of the sum is $`O(K2^{2K})`$. The first term is made of six contributions. Two of them, namely $`2^K(a_0+\lambda (a_1+a_2)+\lambda ^2a_3)`$ and $`2^K(a_0+\lambda (a_4+a_7)+\lambda ^2a_3)`$, are $`O(2^{2K})`$, because of the condition on distances. Among the four remaining contributions, we show how to deal with one of them, the others being handled similarly. This contribution can be written as: $$(a_0+\lambda (a_1+a_4)+\lambda ^2a_5)^K=(a_0+a_1+a_4+a_5)^K\left(1+\frac{(\lambda 1)(a_1+a_4)+(\lambda ^21)a_5}{a_0+a_1+a_4+a_5}\right)^K.$$ (82) We distinguish two cases. Either $`a_0+a_1+a_4+a_51/2`$, and we get trivially: $$(a_0+\lambda (a_1+a_4)+\lambda ^2a_5)^K(a_0+a_1+a_4+a_5)^K=O(2^K),$$ (83) since both terms are $`O(2^K)`$; or $`a_0+a_1+a_4+a_51/2`$, and then: $$\begin{array}{c}\hfill \left\|(a_0+\lambda (a_1+a_4)+\lambda ^2a_5)^K(a_0+a_1+a_4+a_5)^K\right\|\\ \hfill \left\|\left(1+\frac{(\lambda 1)(a_1+a_4)+(\lambda ^21)a_5}{a_0+a_1+a_4+a_5}\right)^K1\right\|=O(K2^K).\end{array}$$ (84) Using again Eq. (78) finishes the proof of Lemma 4.4. $`\mathrm{}`$ ### 4.5 Heuristics for arbitrary $`x`$ For arbitrary $`x`$, the function to minimize in (47) is hard to study analytically. Here we present what we believe to be the correct asymptotic expansion of $`\alpha _{LB}(K,x)`$ at large $`K`$. Hopefully this temptative analysis could be used as a starting point towards a rigorous analytical treatment for any $`x`$. A careful look at the numerics suggests the following Ansatz on the position of the global maximum, at large $`K`$: $$\begin{array}{cc}& a_0=1x+o(1),a_6=x+o(1)\hfill \\ & a_i=o(1)\text{for }i0,6.\hfill \end{array}$$ (85) A second, symmetric, maximum also exists around $`a_0=1x`$, $`a_5=x`$. Plugging this locus into Eq. (47) leads to the following conjecture: ###### Conjecture 4.6 For all $`x(0,1]`$, the asymptotics of $`\alpha _{LB}(x)`$ is given by: $$\underset{K\mathrm{}}{lim}2^K\alpha _{LB}(K,x)=\frac{\mathrm{ln}2+H(x)}{2},$$ (86) and the limit is uniform on any closed sub-interval of $`(0,1]`$. This conjecture is consistent with both our numerical simulations and our result at $`x=\frac{1}{2}`$. ## 5 Proof of Theorem 1.6 Starting with the sharpness criterion for monotone properties of the hypercube given by E. Friedgut and J. Bourgain, we will prove Theorem 1.6 by using techniques and tools developped by N. Creignou and H. Daudé for proving the sharpness of monotone properties in random CSPs. First we make precise some notations for this study on random $`K`$-CNF formula over $`N`$ Boolean variables $`\{x_1,\mathrm{},x_N\}.`$ A $`K`$-clause $`C`$ is given in disjunctive form: $`C=x_1^{\epsilon _1}\mathrm{}x_K^{\epsilon _K}`$ where $`\epsilon _i\{0,1\}`$ ($`x_i^0`$ is the positive literal $`x_i`$ and $`x_i^1`$ is the negative one $`\overline{x_i}`$). A $`K`$-CNF formula $`F`$ is a finite conjunction of $`K`$-clauses, $`\mathrm{\Omega }(F)`$ will denote the set of distinct variables occurring in $`F`$, $`\mathrm{\Omega }(F)\{x_1,\mathrm{},x_N\}`$. In this Boolean framework, $`S(F)`$ the set of satisfying assignments to $`F`$, becomes a subset of $`\{0,1\}^N`$. Now, let us recall how a slight change of our probability measure on formulæ gives a convenient product probability space for studying $`x`$-satisfiability. ### 5.1 $`x`$-unxatisfiability as a monotone property In our case the number of clauses in a random formula $`F_K(N,N\alpha )`$ is fixed to $`M=N\alpha `$. We define another kind of random formula $`G_K(N,N\alpha )`$ by allowing each of the $`𝒩=2^K\left(\genfrac{}{}{0pt}{}{N}{K}\right)`$ possible clauses to be present with probability $`p=\alpha N/𝒩`$. Then, assigning $`1`$ to each clause if it is present and $`0`$ otherwise, the hypercube $`\{0,1\}^𝒩`$ stands for the set of all possible formulæ, endowed with the so-called product measure $`\mu _p`$, where $`p`$ is the probability for $`1`$, and $`1p`$ for $`0`$. More generally, let $`𝒩`$ be a positive integer, a property $`Y\{0,1\}^𝒩`$ is called monotone if , for any $`y,y^{}\{0,1\}^𝒩`$, $`yy^{}`$ and $`yY`$ implies $`y^{}Y`$. In that case $`\mu _p(yY)`$ is an increasing function of $`p[0,1]`$ where $$\mu _p(y_1,\mathrm{},y_𝒩)=p^{\|y\|}(1p)^{𝒩\|y\|}\text{ where }\|y\|=\mathrm{}\{1i𝒩/y_i=1\}.$$ For any non trivial $`Y`$ we can define for every $`\beta ]0,1[`$ the unique $`p_\beta ]0,1[`$ such that: $$\mu _{p_\beta }(yY)=\beta .$$ In our case $`Y`$ will be the property of being $`x`$-unsatisfiable. If we put: $$𝒟=\{(\stackrel{}{\sigma },\stackrel{}{\tau })\{0,1\}^N\times \{0,1\}^Ns.t.d_{\stackrel{}{\sigma }\stackrel{}{\tau }}[Nx\epsilon (N),Nx+\epsilon (N)]\}$$ (87) then $`x`$-unsatisfiability can be read: $$FYS(F)\times S(F)𝒟=\mathrm{}.$$ Observe that the number of clauses in $`G_K(N,N\alpha )`$ is distributed as a binomial law $`\mathrm{Bin}(𝒩,p=\alpha N/𝒩)`$ peaked around its expected value $`p𝒩=\alpha N`$. Therefore, from well known results on monotone property of the hypercube, (JansonLR-99, , page 21 and Corollary 1.16 page 19), our Theorem 1.6 is equivalent to the following result, which establishes the sharpness of the monotone property $`Y`$ under $`\mu _p`$. ###### Theorem 5.1 For each $`K3`$ and $`x,0<x<1`$, there exists a sequence $`\alpha _N(K,x)`$ such that for all $`\eta >0`$: $$\underset{N\mathrm{}}{lim}\mu _p(F\text{ is }xunsatisfiable)=\{\begin{array}{cc}\hfill 1\text{ if }p𝒩=(1\eta )\alpha _N(K,x)N,& \\ \hfill 0\text{ if }p𝒩=(1+\eta )\alpha _N(K,x)N.& \end{array}$$ (88) This theorem will be proved using general results on monotone properties of the hypercube. We state these results below without proof. ### 5.2 General tools The main tool used to prove the existence of a sharp threshold will be a sharpness criterion stemming from Bourgain’s result Friedgut and from a remark by Friedgut on the possibility to strengthen his criterion (Friedgut-05, , Remark following Theorem 2.2). Thus, a slight strengthening of Bourgain’s proof in the appendix of Friedgut combined with an observation made in (CreignouD-03, , Theorem 2.3, page 130) gives the following sharpness criterion: ###### Theorem 5.2 Let $`Y_𝒩\{0,1\}^𝒩`$ be a sequence of monotone properties, then $`Y`$ has a sharp threshold as soon as there exists a sequence $`T_𝒩`$ with $`T_𝒩Y_𝒩`$ such that for any $`\beta ]0,1[`$ and every $`D1`$ the three following conditions are satisfied: $$p_\beta =o(1),$$ (89) $$\mu _{p_\beta }(y\mathrm{s}.\mathrm{t}.zT,zy,\|z\|D)=o(1),$$ (90) $$z_0T,\|z_0\|D\mu _{p_\beta }(yY,yz_0Y\|yz_0)=o(1).$$ (91) We end this subsection by recalling two general results on monotone properties defined on finite sets, established in CreignouD-04 . ###### Lemma 5.3 (CreignouD-04, , Lemma A.1, page 236) Let $`U=\{1,\mathrm{},𝒩\}`$ be partitioned into two sets $`U^{}`$ and $`U^{\prime \prime }`$ with $`\mathrm{\#}U^{}=𝒩^{},\mathrm{\#}U^{\prime \prime }=𝒩^{\prime \prime }`$ and $`𝒩=𝒩^{}+𝒩^{\prime \prime }`$. For any $`uU`$ let us denote $`u^{}=uU^{}`$ and $`u^{\prime \prime }=uU^{\prime \prime }.`$ Let $`Y\{0,1\}^𝒩`$ be a monotone property. For any element $`u`$, let $`𝒜(u)`$ be the set of elements from $`U^{}`$ that are essential for property $`Y`$ at $`u`$: $`𝒜(u)=\left\{iU^{}\text{ s.t. }u\{i\}Y\right\}.`$ Then, for any $`a>0`$ the following holds $$\mu _p(uY,u^{\prime \prime }Y)\frac{1}{(1p)^𝒩^{}}\mu _p(uY,\mathrm{\#}𝒜(u)a)+\frac{ap}{(1p)^𝒩^{}}.$$ For the second result we consider a sequence of monotone properties $`Y_𝒩\{0,1\}^𝒩`$. For any fixed $`u\{0,1\}^𝒩`$, $`_j(u)`$ will be the set of collections of $`j`$ elements such that one can reach property $`Y`$ from $`u`$ by adding this collection, thus $`\mathrm{\#}_j(u)\left(\genfrac{}{}{0pt}{}{𝒩}{j}\right)`$. ###### Lemma 5.4 (CreignouD-04, , Lemma A.2, page 237) Let $`Y_𝒩\{0,1\}^𝒩`$ be a sequence of monotone properties. For any integer $`j1`$, for any $`b>0`$ and as soon as $`𝒩p`$ tends to infinity, the following estimate holds $$\mu _p\left(uY,\mathrm{\#}_j(u)b\left(\genfrac{}{}{0pt}{}{𝒩}{j}\right)\right)=o(1),$$ $$_j(u)=\{\{i_1,\mathrm{},i_j\},1i_1<\mathrm{}<i_j𝒩,\text{ such that }u\{i_1,\mathrm{},i_j\}Y\}.$$ ### 5.3 Proof of Theorem 5.1 (main steps) As usual, the first two conditions $`(\text{89})`$ and $`(\text{90})`$ are easy to verify for the $`x`$-unsatisfiability property. For the first one we have: $$\mu _p(F\text{ is }x\text{-satisfiable})\mu _p(F\text{ is satisfiable})2^N(1p)^{\left(\genfrac{}{}{0pt}{}{N}{K}\right)}.$$ This shows that $`p_\beta {\displaystyle \frac{N\mathrm{ln}(2)\mathrm{ln}(1\beta )}{\left(\genfrac{}{}{0pt}{}{N}{K}\right)}}`$ , thus for $`x`$-unsatisfiability we get: $$\beta ]0,1[p_\beta (N)=O(N^{1K}).$$ (92) For the second condition, let $`H(F)`$ be the $`K`$-uniform hypergraph associated to a formula $`F`$: its vertices are the $`\mathrm{\Omega }(F)`$ variables occurring in $`F`$, each index set of a clause $`C`$ in $`F`$ corresponds to an hyperedge. Let us recall, see KL-02 , that a $`K`$-uniform connected hypergraph with $`v`$ vertices and $`w`$ edges is called a hypertree when $`(K1)wv=1`$; it is said to be unicyclic when $`(K1)wv=0`$, and complex when $`(K1)wv1`$. Let $`T`$ be the set of formulæ $`F`$ such that $`H(F)`$ has at least one complex component. We will rule out $`(\text{90})`$ (and also $`(\text{91})`$) by using the following result on non complex formulæ, the proof of which is deferred to the next subsection: ###### Lemma 5.5 Let $`K3`$. If $`G`$ is a $`K`$-CNF-formula on $`v`$ variables whose associated hypergraph is an hypertree or unicyclic then for all integer $`d\{0,\mathrm{},v\}`$ there exits $`(\stackrel{}{\sigma },\stackrel{}{\tau })S(G)\times S(G)`$ such that $`d_{\stackrel{}{\sigma }\stackrel{}{\tau }}=d.`$ In particular, this result shows that any $`x`$-unsatisfiable formula has at least one complex component, i. e. $`TY.`$ Then observe that there is $`O(N^{(K1)s1})`$ distinct complex components of size $`s`$ with $`N`$ vertices. Thus we get for all $`p:`$ $`\mu _p(F\mathrm{s}.\mathrm{t}.GT,GF,\|G\|D){\displaystyle \underset{sD}{}}O(N^{(K1)s1})p^s,`$ and $`(\text{90})`$ follows from $`(\text{92})`$ In order to prove $`(\text{91})`$, let us introduce some tools inspired of CreignouD-04 . For each positive integer $`t`$ and $`\mathrm{\Delta }=(\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_t)\{0,1\}^t`$, a $`\mathrm{\Delta }`$-assignment is an assignment for which the $`t`$ first values of the variables are equal to $`\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_t`$. Then $`S_\mathrm{\Delta }(F)`$ will denote the set of satisfying $`\mathrm{\Delta }`$-assignments to $`F`$: $`S_\mathrm{\Delta }(F)S(F)\{0,1\}^N`$. For any pair of $`t`$-tuples $`(\mathrm{\Delta },\mathrm{\Delta }^{})\{0,1\}^t\times \{0,1\}^t`$ we define $`Y^{\mathrm{\Delta },\mathrm{\Delta }^{}}`$: $$FY^{\mathrm{\Delta },\mathrm{\Delta }^{}}S_\mathrm{\Delta }(F)\times S_\mathrm{\Delta }^{}(F)𝒟_x=\mathrm{}.$$ Observe that $`Y^{\mathrm{\Delta },\mathrm{\Delta }^{}}`$ is a monotone property containing $`Y`$. Now we come back to $`(\text{91})`$ with $`F_0T`$, so that the hypergraph associated to the booster formula $`F_0`$ has no complex components. $`S(F_0)\mathrm{}`$ and w.l.o.g. we can suppose that $`\mathrm{\Omega }(F_0)=\{1,\mathrm{},t\}`$. Then, for $`FY`$ such that $`FF_0`$ with $`FF_0Y`$, let $`F^{\prime \prime }`$ denote the largest subformula of $`F`$ such that $`\mathrm{\Omega }(F^{\prime \prime })\{1,\mathrm{},t\}=\mathrm{}`$. We have the two following claims whose proof is postponed to the next subsection. ###### Claim 5.6 For any $`(\mathrm{\Delta },\mathrm{\Delta }^{})S(F_0)\times S(F_0),FF_0Y^{\mathrm{\Delta },\mathrm{\Delta }^{}}`$. ###### Claim 5.7 There exits $`(\mathrm{\Delta },\mathrm{\Delta }^{})S(F_0)\times S(F_0)`$ such that $`F^{\prime \prime }Y^{\mathrm{\Delta },\mathrm{\Delta }^{}}`$. Thus $`(\text{91})`$ is proved as soon as for any $`\beta ]0,1[`$ and $`(\mathrm{\Delta },\mathrm{\Delta }^{})\{0,1\}^t\times \{0,1\}^t`$: $$\mu _{p_\beta }(FF_0Y^{\mathrm{\Delta },\mathrm{\Delta }^{}},F^{\prime \prime }Y^{\mathrm{\Delta },\mathrm{\Delta }^{}}\|FF_0)=o(1).$$ (93) The two first events in the R.H.S. of (93) do not depend on the set of clauses in $`F_0`$ thus by independence under the product measure and recalling that $`Y^{\mathrm{\Delta },\mathrm{\Delta }^{}}`$ is a monotone property we are led to prove that: $$\mu _{p_\beta }(FY^{\mathrm{\Delta },\mathrm{\Delta }^{}},F^{\prime \prime }Y^{\mathrm{\Delta },\mathrm{\Delta }^{}})=o(1).$$ From $`(\text{92})`$ we know that $`p_\beta (N)=O(N^{1K})`$. Let $`𝒩^{}=\mathrm{\Theta }(N^{K1})`$ be the number of clauses having at least one variable in $`\{1,\mathrm{},t\}`$, then Lemma 5.3, applied to the monotone property $`Y^{\mathrm{\Delta },\mathrm{\Delta }^{}}`$, shows that the above assertion is true as soon as we are able to prove that for all $`\gamma >0`$: $$\mu _{p_\beta }(FY^{\mathrm{\Delta },\mathrm{\Delta }^{}},\mathrm{\#}𝒜_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)\gamma N^{K1})=o(1).$$ (94) where $`𝒜_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$ is the set of $`K`$-clauses $`C`$ on $`N`$ variables having at least one variable in $`\{x_1,\mathrm{},x_t\}`$ and such that $`FCY^{\mathrm{\Delta },\mathrm{\Delta }^{}}`$. Then let $`_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$ be the set of collections of $`(K1)`$ $`K`$-clauses $`\{C_1,\mathrm{},C_{K1}\}`$ such that $`FC_1\mathrm{}C_{K1}Y^{\mathrm{\Delta },\mathrm{\Delta }^{}}`$. From lemma 5.3 we deduce that (94) is true as soon as the following result is proved: ###### Lemma 5.8 For all $`t,K3,\gamma >0`$ and $`(\mathrm{\Delta },\mathrm{\Delta }^{})\{0,1\}^t\times \{0,1\}^t`$, there exits $`\theta >0`$ such that for all $`N`$, the following holds: $$\mathrm{\#}𝒜_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)\gamma N^{K1}\mathrm{\#}_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)\theta N^{K(K1)}.$$ (95) Again the proof of this last result is deferred to the next subsection that furnishes a detailed and complete proof of Theorem 5.1. ### 5.4 Detailed proofs #### 5.4.1 Lemma 5.5 Proof: When $`G`$ has a leaf-clause, that is a clause $`C=x_1^{\epsilon _1}\mathrm{}x_K^{\epsilon _K}`$ having only one variable, say $`x_1`$, in common with $`GC`$, the assertion can be proved by induction on the number of clauses in $`G`$. Indeed from a pair of satisfying assignments $`(\stackrel{}{\sigma },\stackrel{}{\tau })S(GC)\times S(GC)`$ with $`d_{\stackrel{}{\sigma }\stackrel{}{\tau }}=d`$ and a pair of satisfying assignments at distance $`d^{}\{0,\mathrm{},K1\}`$ for $`C^{}=x_2^{\epsilon _2}\mathrm{}x_K^{\epsilon _K}`$, one gets a pair of satisfying assignments at distance $`d+d^{}`$. But $`C^{}`$ is a $`K1`$-clause, thus for any $`d^{}\{0,\mathrm{},K1\}`$ $`C^{}`$ has a pair of satisfying assignments at distance $`d^{}`$. When any $`K`$-clause $`C_i`$ of $`G=C_1\mathrm{}C_l`$ has exactly two variables in common with $`GC_i`$ then we can write $`C_1=x_1^{\mu _1}x_2^{\nu _2}C_1^{},C_2=x_2^{\mu _2}x_3^{\nu _3}C_2^{},\mathrm{},C_l=x_l^{\mu _l}x_1^{\nu _1}C_l^{}`$ where the $`C_j^{}`$ are $`(K2)`$-clauses. A variable in $`C_j^{}`$ occurs exactly once in formula $`G`$ and the set of variables in these $`C_j^{}`$ is equal to $`\{x_{l+1},\mathrm{},x_v\}`$. In particular this set is disjoint from the set of variables of the $`2`$-CNF formula $`(x_1^{\mu _1}x_2^{\nu _2})(x_2^{\mu _2}x_3^{\nu _3})\mathrm{}(x_l^{\mu _l}x_1^{\nu _1})`$. First observe that this $`2`$-CNF cyclic formula has always a satisfying assignment $`(\sigma _1,\mathrm{},\sigma _l)`$ and together with any truth value for the $`(x_j,j>l)`$ it gives a satisfying assignment for $`G`$. Thus, for $`G`$, one gets a pair of satisfying assignments at distance $`d`$ for any $`dvl`$. Second, as $`\mathrm{\Omega }(C_j^{})\mathrm{\Omega }(C_k^{})=\mathrm{}`$ when $`jk`$ a satisfying assignment $`\sigma _{l+1},\mathrm{},\sigma _v`$ can easily be found for $`C_1^{}\mathrm{}C_l^{}`$. Together with any truth values of the $`(x_i,il)`$ it gives a satisfying assignment for $`G`$. Then, from the satisfying assignment $`(\sigma _1,\mathrm{},\sigma _l,1\sigma _{l+1},\mathrm{},1\sigma _v)`$ one gets, for any $`dvl`$, a pair of satisfying assignments at distance $`d`$. $`\mathrm{}`$ #### 5.4.2 Claims 5.6 and 5.7 Proof: Observe that any SAT-$`x`$-pair $`(\stackrel{}{\sigma },\stackrel{}{\tau })`$ for $`FF_0`$ with $`(\sigma _1,\mathrm{},\sigma _t)S(F_0)`$ and $`(\tau _1,\mathrm{},\tau _t)S(F_0)`$ is also a SAT-$`x`$-pair for $`F`$. This proves the first claim by contradiction. For the second claim, $`FF_0Y`$ so there exits a SAT-$`x`$-pair $`(\stackrel{}{\sigma },\stackrel{}{\tau })S(FF_0)\times S(FF_0)`$. By construction, the set of satisfying assignment of $`F^{\prime \prime }`$ does not depend on the first $`t`$ coordinates. Let $`d_t`$ be the Hamming distance between $`(\sigma _1,\mathrm{}\sigma _t)`$ and $`(\tau _1,\mathrm{}\tau _t)`$. We know that all components of the hypergraph associated to formula $`F_0`$ are simple and lemma $`(\text{5.5})`$ shows that there exits $`(\sigma _1^{},\mathrm{}\sigma _t^{})S(F_0)`$ and $`(\tau _1^{},\mathrm{}\tau _t^{})S(F_0)`$ such that $`d_{\stackrel{}{\sigma }^{}\stackrel{}{\tau }^{}}=d_t`$. Hence $`(\sigma _1^{},\mathrm{}\sigma _t^{},\sigma _{t+1},\mathrm{},\sigma _N)`$ and $`(\tau _1^{},\mathrm{}\tau _t^{},\tau _{t+1},\mathrm{},\tau _N)`$ form now a SAT-$`x`$-pair for $`F^{\prime \prime }`$, thus proving the second claim. $`\mathrm{}`$ #### 5.4.3 Lemma 5.8 Proof: In ErdosS-82 , Erdös and Simonovits proved that any sufficiently dense uniform hypergraph always contains specific subhypergraphs. In particular they considered a generalization of the complete bipartite graph specified by two integers $`h2`$ and $`m1`$. Let us denote by $`K_h(m)`$ the $`h`$-uniform hypergraph with $`hm`$ vertices partitioned into $`h`$ classes $`V_1,\mathrm{},V_h`$ with $`\mathrm{\#}V_i=m`$ and whose hyperedges are those $`h`$-tuples, which have exactly one vertex in each $`V_i`$. Thus $`K_h(m)`$ has $`m^h`$ hyperedges, for $`h=2`$ it is a complete bipartite graph $`K(m,m)`$. For proving Lemma 5.8, we need a small variation on a result of Erdös and Simonovits which differs only in that it deals with ordered $`h`$-tuples as opposed to sets of size $`h`$. More precisely, let us consider hypergraphs on $`n`$ vertices, say $`\{x_1,\mathrm{},x_n\}`$, we will say that two disjoint subsets of vertices $`A`$ and $`B`$ verify $`A<B`$ if for all $`x_i`$ in $`A`$ and all $`x_j`$ in $`B`$ we have $`i<j`$. Let $`H`$ be an $`h`$-uniform hypergraph with vertex set $`\{x_1,\mathrm{},x_n\}`$, then any $`h`$-uniform subhypergraph $`K_h(m)`$ with $`V_1<\mathrm{}<V_h`$ is called an ordered copy of $`K_h(m)`$ in $`H`$. Thus, the ordered version of the theorem from Erdös and Simonovits about supersaturated uniform hypergraphs (ErdosS-82, , Corollary 2, page 184) can be stated as follows. ###### Theorem 5.9 (Ordered Erdös-Simonovits) Given $`c>0`$ and two integers $`h2`$ and $`m1`$, there exist $`c^{}>0`$ and $`N`$ such that for all integers $`nN`$, if $`H`$ is a $`h`$-uniform hypergraph over $`n`$ vertices having at least $`c\left(\genfrac{}{}{0pt}{}{n}{h}\right)`$ hyperedges then $`H`$ contains at least $`c^{}n^{hm}`$ ordered copies of $`K_h(m)`$. We will also use the following observation made when one consider an assignment of two colours, say $`0`$ and $`1`$, to the hyperedges of $`K_h(m)`$. First let’s say that a vertex $`s`$ is $`c`$-marked if $`s`$ belongs to at least one $`c`$-colored hyperedge. A subset of vertices $`S`$ is said $`c`$-marked if any $`s`$ in $`S`$ is $`c`$-marked. ###### Claim 5.10 Let $`h2`$, $`m1`$, and $`V_1,\mathrm{},V_h`$ the partition associated to $`K_h(m)`$. Consider an assignment of two colours to the $`m^h`$ hyperedges of $`K_h(m)`$, then at least one of the $`V_i`$ is marked. Indeed, suppose that $`V_1,\mathrm{},V_h`$ are not $`c`$-marked. Now consider a vertex $`sV_1`$ then $`s`$ is $`(1c)`$ marked else by construction of $`K_h(m)`$, $`V_i`$ would be $`c`$-marked for all $`i2`$. Hence $`V_1`$ becomes $`(1c)`$-marked. Now let us show (95), in other words that for any $`K`$-CNF formula $`F`$ such that $`𝒜_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$ is dense then $`_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$ is also dense. For more readability we will restrict our attention to the special case $`K=3`$, in using the above fact the proof will be easily extendable to any $`K3`$. Suppose there exist $`\mathrm{\Theta }(N^2)`$ clauses in $`𝒜_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$ then, by the pigeon hole principle, at least for one of the eight types of clause we can find $`\mathrm{\Theta }(N^2)`$ clauses of this type in $`𝒜_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$. Suppose, for example, that $$\mathrm{\#}\{C=\overline{x_{i_1}}x_{i_2}\overline{x_{i_3}},1i_1<i_2<i_3N,i_1t,FCY^{\mathrm{\Delta },\mathrm{\Delta }^{}}\}=\mathrm{\Theta }(N^2).$$ From well chosen elements in $`𝒜_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$ we now exhibit an element in $`_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$. We consider the graph $`H(F)`$ associated to formula $`F`$: the set of vertices is $`\{1,\mathrm{},N\}`$ and for each $`C=\overline{x_{i_1}}x_{i_2}\overline{x_{i_3}}𝒜_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$ we create an edge $`\{i_2,i_3\}`$. Let $`(\stackrel{}{\sigma },\stackrel{}{\tau })`$ be a SAT-$`x`$-pair for $`F`$, then either $`\sigma S(C)`$ or $`\tau S(C)`$. Now, following a fixed ordering on the set of pairs of thruth assignments we put the colour $`0`$ on the non colored edge $`\{i_2,i_3\}`$ if $`\sigma _{i_2}=0`$ and $`\sigma _{i_3}=1`$ else we put the color $`1`$, having in this case $`\tau _{i_2}=0`$ and $`\tau _{i_3}=1`$. Now, let’s take an ordered copy of $`K(3,3)`$ in $`H(F)`$ with partition $`A=\{j_1,j_2,j_3\}`$ and $`B=\{j_4,j_5,j_6\}`$. From Fact 5.10 we know that one part, say $`A`$, is marked. In such a case we have $`\sigma _{j_1}=0,\sigma _{j_2}=0,\sigma _{j_3}=0`$ (A is $`0`$-marked) or $`\tau _{j_1}=0,\tau _{j_2}=0,\tau _{j_3}=0`$ (A is $`1`$-marked) hence $`(\stackrel{}{\sigma },\stackrel{}{\tau })`$ is no longer a SAT-$`x`$-pair for $`F(x_{j_1}x_{j_2}x_{j_3})`$. If $`B`$ is marked then $`(\stackrel{}{\sigma },\stackrel{}{\tau })`$ is no longer a SAT-$`x`$-pair for $`F(\overline{x_{j_4}}\overline{x_{j_5}}\overline{x_{j_6}})`$. Thus in any case $`\{(x_{j_1}x_{j_2}x_{j_3}),(\overline{x_{j_4}}\overline{x_{j_5}}\overline{x_{j_6}})\}_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$. By hypothesis $`H(F)`$ is a dense graph so from Theorem 5.9 we can find $`\mathrm{\Theta }(N^6)`$ copies of $`K(3,3)`$ in $`H(F)`$. The above construction provide $`\mathrm{\Theta }(N^6)`$ elements in $`_{\mathrm{\Delta },\mathrm{\Delta }^{}}(F)`$ thus proving that this set is also dense. $`\mathrm{}`$ ### 5.5 A general sharpness result Note that the above proof does not use any information about the shape of the set $`𝒟`$ defining the $`x`$-unsatisfiability in terms of a subset of $`\{0,\mathrm{},N\}`$, namely the interval $`[Nx\epsilon (N),Nx+\epsilon (N)]`$ (see $`(\text{87})`$). Actually we can consider properties defined by a non empty proper subset of $`\{0,\mathrm{},N\}`$ and we have proved the following general result: ###### Theorem 5.11 Let $`J_N`$ be a non empty subset of $`\{0,\mathrm{},N\}`$ and consider $$𝒟_J=\{(\stackrel{}{\sigma },\stackrel{}{\tau })\{0,1\}^N\times \{0,1\}^Ns.t.d_{\stackrel{}{\sigma }\stackrel{}{\tau }}J_N\}.$$ Let $`K3`$ and $`Y_J`$ be the set of $`K`$-CNF formula defined as: $$FY_JS(F)\times S(F)𝒟_J=\mathrm{}.$$ Then, $`Y_J`$ is a monotone property exhibiting a sharp threshlold. On one hand, any upper bound for the satisfiability threshold, for instance $`(\text{92})`$, is an upper bound for all $`Y_J`$ threshold. On the other hand, lemma 5.5 tells us that a non complex formula does not belongs to $`Y_J`$. Then, from KL-02 , we know that w.h.p a formula whose ratio between the number of clauses and the number of varibles is less than $`1/K(K1)`$, has no complex component. Thus it provides a lower bound for all $`Y_J`$ threshold. ## 6 Discussion and Conclusion We have developed a simple and rigorous probabilistic method which is a first step towards a complete characterization of the clustered hard-SAT phase in the random satisfiability problem. Our result is consistent with the clustering picture and supports the validity of the one-step replica symmetry breaking scheme of the cavity method for $`K8`$. The study of $`x`$-satisfiability has the advantage that it does not rely on a precise definition of clusters. Indeed, it is important to stress that the “appropriate” definition for clusters may vary according to the problem at hand. The natural choice seems to be the connected components of the space of SAT-assignments, where two adjacent assignments have by definition Hamming distance $`1`$. However, although this naive definition seems to work well on the satisfiability problem, it raises major difficulties on some other problems. For instance, in $`q`$-colorability, it is useful to permit color exchanges between two adjacent vertices in addition to single-vertex color changes. In XORSAT, the naive definition is inadequate, since jumps from solution to solution can involve a large, yet finite, Hamming distance due to the hard nature of linear Boolean constraints MontanariSemerjian05-2 . On the other hand, the existence of a gap in the $`x`$-satisfiability property is stronger than the original clustering hypothesis. Clusters are expected to have a typical size, and to be separated by a typical distance. However, even for typical formulas, there exist atypical clusters, the sizes and separations of which may differ from their typical values. Because of this variety of cluster sizes and separations, a large range of distances is available to pairs of SAT-assignments, which our $`x`$-satisfiability analysis takes into account. What we have shown suggests that, for typical formulas, the maximum size of all clusters is smaller than the minimum distance between two clusters (for a certain range of $`\alpha `$ and $`K8`$). This is a sufficient condition for clustering, but by no means a necessary one. As a matter of fact, our large $`K`$ analysis conjectures that $`\alpha _1(K)`$ (the smaller $`\alpha `$ such that Conjecture 1.4 is verified) scales as $`2^{K1}\mathrm{ln}2`$, whereas $`\alpha _d(K)`$ (where the replica symmetry breaking occurs) and $`\alpha _s(K)`$ (where the one-step RSB Ansatz is supposed to be valid) scale as $`2^K\mathrm{ln}K/K`$ MMZ-RSA . According to the physics interpretation, in the range $`\alpha _s(K)<a<\alpha _1(K)`$, there exist clusters, but they are not detected by the $`x`$-satisfiability approach. This limitation might account for the failure of our method for small values of $`K`$ — even though more sophisticated techniques for evaluating the $`x`$-satisfiability threshold $`\alpha _c(K,x)`$ might yield some results for $`K<8`$. Still, the conceptual simplicity of our method makes it a useful tool for proving similar phenomena in other systems of computational or physical interest. A better understanding of the structure of the space of SAT-assignments could be gained by computing the average configurational entropy of pairs of clusters at fixed distance, which contains details about how intra-cluster sizes and inter-cluster distances are distributed. This would yield the value of the $`x`$-satisfiability threshold. Such a computation was carried out at a heuristic level within the framework of the cavity method for the random XORSAT problem MoraMezard06 , and should be extendable to the satisfiability problem or to other CSPs. This work has been supported in part by the EC through the network MTR 2002-00319 ‘STIPCO’ and the FP6 IST consortium ‘EVERGROW’. This paper, signed in alphabetic order, is based on previous work by Mora Mézard and Zecchina reported in Sec. 1–4, 6. The proof in Sec. 5 is due to Daudé.
warning/0506/astro-ph0506258.html	ar5iv	text	# Chandra Study of X-Ray Point Sources in the Early-Type Galaxy NGC 4552 (M89) ## 1 INTRODUCTION Elliptical and S0 galaxies are luminous sources in the X-ray sky (e.g., Forman et al. 1985). In terms of their X-ray-to-optical luminosity ratios, these early-type galaxies can generally be divided into two categories: X-ray bright galaxies and X-ray faint galaxies. Both X-ray faint and X-ray bright galaxies reveal a hard X-ray component whose intensity is roughly proportional to the optical luminosity of the galaxy. In X-ray bright early-type galaxies, the hard spectral component is often overwhelmed by the emission of the hot diffuse gas. However the hard X-ray component appears to be significant in X-ray faint galaxies and, in some cases, dominates the X-ray spectrum. Trinchieri and Fabbiano (1985) and other later studies of the non-thermal emission ascribed it to the contribution of a large number of low mass X-ray binaries (LMXBs), such as observed in the bulges of M31 and our Galaxy (e.g., White et al. 1995). With the advent of the high spatial resolution observations of the Chandra X-ray observatory, hundreds of point sources have been resolved in a number of nearby elliptical and S0 galaxies and the LMXB-origin has been confirmed (e.g., Sarazin et al. 2000). Most of the resolved X-ray point sources in early-type galaxies are likely LMXBs and they may be useful in deciphering the evolutionary history of their host galaxies and stars therein, and constraining our understanding of the physics of compact stars. In the past few years Chandra observations have revealed that the X-ray luminosities of the resolved point sources in early type galaxies span a wide range from the typical observational limit of a few $`10^{37}`$ erg s<sup>-1</sup> to over $`10^{39}`$ erg s<sup>-1</sup> (e.g., Sarazin et al. 2000; Angelini et al. 2001; Kundu et al. 2002). The observed XLFs show quite similar, but not exactly the same profiles among galaxies. For example, in NGC 4697, Sarazin et al. (2000) reported that there is a knee on the XLF at around the Eddington luminosity of normal accreting neutron stars with a mass $`1.4M_{}`$ ($`L_{\mathrm{Edd}}2\times 10^{38}`$ ergs s<sup>-1</sup>), and thus suggested that the break might be an universal feature that can be used as a distance indicator. The existence of a similar break was confirmed by Kundu et al. (2002) in NGC 4472 and by Blanton et al. (2001) in NGC 1553. However, by analyzing a sample of 14 E/S0 galaxies, Kim and Fabbiano (2004) argued that the position of the break is significantly higher than $`L_{\mathrm{Edd}}`$ of normal neutron stars. In NGC 720, Jeltema et al. (2003) showed an even higher break at $`L_\mathrm{b}=2.10_{0.2}^{+0.2}\times 10^{39}`$ ergs s<sup>-1</sup> (although this is based on a distance derived from an adopted $`H_0`$ of 50 km s<sup>-1</sup> Mpc<sup>-1</sup>). Similar breaks or upper cutoffs at substantially high luminosities have been found by Sivakoff et al. (2003) in NGC 4365 and NGC 4382, and by Jord$`\stackrel{´}{\mathrm{a}}`$n et al. (2004) in M87, M49 and NGC 4697. Moreover, in NGC 1600 Sivakoff et al. (2004) found that the break is not needed to fit the observed XLF. The obvious disagreements between these works cast doubts on the feasibility of using $`L_\mathrm{b}`$ as a reliable standard candle for distance determinations. On the other hand, it has been revealed that in early-type galaxies about 4% of the GCs host bright LMXBs, and about 18–70% of the resolved LMXBs are found to be associated with GCs with a preference on the optically brighter and redder ones (e.g., Kundu et al. 2002; Sarazin et al. 2000, 2001, 2003). In most cases, such a correlation with color is attributed to metallicity (Kundu et al. 2003). These results imply that GCs are an efficient breeding ground for LMXBs, as it is true in our own Galaxy (Katz 1975; Clark 1975). The origin of the LMXBs that are not associated with GCs (field LMXBs) is still unclear. It has been speculated in many published papers (e.g., Kundu et al. 2002, Maccarone et al 2003) that the majority of them may be an ejected GC population, or have been left in the field after GC disaggregation. Still, the possibility that they are a true field population and were actually formed in situ cannot be ruled out at present (Maccarone et al. 2003). The comparisons between the X-ray properties of the GC LMXBs and field LMXBs should help answer the question. As of today, it appears that in some cases the difference between the spectra of GC LMXBs and field LMXBs is almost undistinguishable in the X-ray band, strongly supporting a similar origin. The only exception may be NGC 4472, for which Maccarone et al. (2003) showed that the GC LMXBs tend to be slightly harder. In this paper, we present a Chandra study of the X-ray point source population in the X-ray bright galaxy NGC 4552 (M89). NGC 4552 is located at $`z=0.001134`$ (Smith et al. 2000) in the Virgo cluster. It is usually classified as an elliptical galaxy, although in some cases (e.g., Ferrari et al. 1999) it is considered a S0. It is also a bright radio source, with a strong compact core and a relatively flat spectrum (Filho et al. 2000; Wrobel & Heeschen 1984). HST observations show that from 1991 to 1996 the intensity of the central, unresolved source of this galaxy changed by a factor of several in the near UV band, along with the appearance of some UV/optical emission lines (Renzini et al. 1995; Cappellari et al. 1999). This indicates that NGC 4552 harbors a mini-AGN at the center, as has been inferred in the X-ray band by Colbert and Mushotzky (1999) who studied the nuclear source with ROSAT and ASCA. Throughout this paper, by assuming the Hubble constant to be $`H_0`$ = 70 km s<sup>-1</sup> Mpc<sup>-1</sup> (O’Sullivan et al. 2001 and references therein) we adopt a distance to NGC 4552 of 17.1 Mpc, which is slightly larger than that derived by infrared surface brightness fluctuation analysis ($`15.4\pm 1.0`$ Mpc; Tonry et al. 2001). We quote errors at the 90% confidence level unless mentioned otherwise. ## 2 OBSERVATION AND DATA REDUCTION NGC 4552 was observed on April 22–23, 2001 with the CCD 2, 3, 6, 7 and 8 of the Chandra Advanced CCD Imaging Spectrometer (ACIS) for a total exposure of 56.8 ks. The center of the galaxy was positioned on the ACIS S3 chip (CCD 7) with an offset of $`0.64^{}`$ from the nominal pointing for the S3 chip, so the entire galaxy was covered by the S3 chip. The CCD temperature was $`120^{}`$C. The events were telemetried in the Very Faint mode, and the data were collected with frame times of 3.1 s. In the analysis that follows, we used the CIAO 2.3 software to process the data acquired from the S3 chip only. In order to use the latest calibration, we started with the Level-1 data. We only kept events with ASCA grades 0, 2, 3, 4, and 6, and removed bad pixels, bad columns, and columns adjacent to bad columns and node boundaries. In order to identify occasional intervals of high background (“background flares”), whose effects are particularly significant on the backside-illuminated S1 and S3 chips, we extracted and examined lightcurves of the background regions on the S3 chip in 2.5–7 keV where the background flares are expected to be most visible. We found that the detected intervals contaminated by the particle events that raise the count rate to over 20% more than its mean value is less than 9% of the total exposure. We have excluded these intervals and used a net exposure time of 52.2 ks in our analysis. Also we corrected the aspect offset of the observation ($`\delta `$R.A. = $`0.17^{\prime \prime }`$ and $`\delta `$Dec. = $`0.12^{\prime \prime }`$). We limited the spectral analysis to the $`0.77`$ keV energy band in order to avoid the effects of calibration uncertainties at lower energies and instrumental background at higher energies. We extracted all of the spectra in the pulse height-invariant (PI) channels, and performed model fittings with XSPEC v11.2.0. Because there has been a continuous degradation in the ACIS quantum efficiency (QE) since launch, we applied ACISABS to correct the created ARF files before using them to fit the spectra. ## 3 X-RAY IMAGE In Figure 1, we show the raw Chandra image of NGC 4552 in 0.3–10 keV, which has not been corrected for either exposure or background. For a comparison, we overlap the DSS optical intensity contours on the X-ray image. It can be clearly seen that the spatial distribution of the X-ray emissions is nearly symmetric within about 9<sup>′′</sup> (0.7 kpc). Outside the nuclear region the distribution of the X-rays is elongated roughly in the north-south direction out to about 50<sup>′′</sup> (4.1 kpc), where it does not follow the profile of the optical light. In the region about $`20^{\prime \prime }`$ east of the center there is a lack of diffuse X-ray emission. In Figure 2, we show the smoothed X-ray image with a minimum significance of 3 and a maximum significance of 5. In the figure we use a large circle to indicate the region within the 4 effective radii (4 $`R_\mathrm{e}`$), where 1 $`R_\mathrm{e}=0.49^{}`$ or 2.4 kpc (de Vaucouleurs et al. 1992). The locations of the detected X-ray point sources are also marked, which show a clear tendency to concentrate toward the center. The details of the detection and analysis of these point sources are presented in §4. There is a lack of diffuse X-rays in the region about $`20^{\prime \prime }`$ east of the center, while the region about $`36^{\prime \prime }`$ south-west of the center lacks bright X-ray point sources. We find that the Chandra X-ray position of the central point source (Src 1; see §4.6), the brightest one in the field as well as the peak of the diffuse X-ray emission, coincides with the optical/IR center of the galaxy (Monet et al. 1998; Cutri et al. 2003) within $`0.5^{\prime \prime }`$. This is the X-ray counterpart of the mini-AGN identified through its UV/optical activity (Cappellari et al. 1999). We do notice that by using the ROSAT HRI data Colbert and Mushotzky (1999) reported an offset of $`4.6^{\prime \prime }`$ between the position of the compact X-ray source in the nuclear region and the position of the optical photometric center of the galaxy. But since this offset is much smaller than the uncertainty of the ROSAT HRI positioning, which is typically $`10^{\prime \prime }`$, our results are not in conflict. Moreover, the Chandra X-ray position of the central source is also in good agreement with the position of the central radio source (Nagar et al. 2002) to within about $`1^{\prime \prime }`$. ## 4 X-RAY POINT SOURCES ### 4.1 Detections We detected X-ray point sources on the ACIS S3 image using the CIAO tool celldetect. The default signal-to-noise threshold for source detection was set to be 3, and the energy range used for detection was restricted to 0.3–10 keV for better statistics. We have cross-checked the results both by using the CIAO tool wavdetect and by eye in either the 0.3–10 keV or 0.7–7 keV images. We detected a total of 79 sources that exceeded the detection threshold. Of these, 47 lie within the 4 $`R_\mathrm{e}`$ region. We estimate that the minimum detection for a point source is approximately $`3.13\times 10^4`$ counts s<sup>-1</sup> (16 counts for 52.2 ks) at $`r4R__\mathrm{e}`$. Notice that this value is higher than that in Sivakoff et al. (2003) and Sarazin et al. (2001) by about 45% and 33%, respectively. By reducing the signal-to-noise threshold down to 2.7, 5 more sources can be detected within the 4 $`R_\mathrm{e}`$ region. Based on visual inspection two of them are clearly fake because the photon distribution within the detection cell does not follow that of a point source. Considering the large uncertainties due to the sample incompleteness for very faint sources, in this paper we conservatively adopt a signal-to-noise threshold of 3 and focus our study on the 47 sources within the 4 $`R_\mathrm{e}`$ region. We list the properties of all 47 sources in Table 1, where we sort them in the order of increasing projected distance d from the center of the galaxy. We arrange the columns as follows: (1) source number; (2)-(3): right ascension and declination (J2000); (4) projected distance d from the center of the galaxy; (5) count rate and its error; (6) significance of the detection; (7) intrinsic X-ray luminosity $`L_X`$, assuming the source is located at the distance of NGC 4552 and only subjected to the Galactic absorption ($`2.56\times 10^{20}`$ cm<sup>-2</sup>; Dickey & Lockman, 1990; see §4.5); (8)-(9): hardness ratios (see §4.3); and (10) notes. ### 4.2 Variability of Sources We extracted the lightcurves for each of the 47 point sources over the duration of the Chandra observation, excluding the intervals of strong background flares. The extractions were made in 0.3–10 keV for better statistics, since we found that after correcting for the background the results obtained in this energy band are consistent with those obtained in 0.7–7 keV. We calculated the Kolmogoroff-Smirov (KS) statistic for each of the lightcurves against the null hypothesis that the count rate of the source plus the background is uniform over the effective exposure time. If the count rate is temporally invariant, the cumulative fraction of the count is expected to be a diagonal from 0 to 1. We find that two sources, Src 15 (174 counts) and Src 28 (107 counts), have a less than 5% probability of being invariable. In order to examine if these results are caused by local background fluctuations, we also applied the KS test to the lightcurves extracted from the background regions adjacent to the two sources. In neither case can we find any significant temporal variability in the background. Therefore, we conclude that both sources are intrinsically variable. Assuming that they are at the distance of NGC 4552 and using the best-fit absorbed power-law spectral model for all the resolved sources (§4.4), we estimate that Src 15 and Src 28 have luminosities of $`1.17_{0.30}^{+0.11}\times 10^{39}`$ ergs s<sup>-1</sup> and $`5.12_{1.01}^{+1.00}\times 10^{38}`$ ergs s<sup>-1</sup>, respectively. In Figure 3, we present their background-corrected lightcurves that show clear variations on $`1.52`$ hr timescales, together with the lightcurve of the central source (Src 1), which is the brightest one in the field. In terms of the KS test, the temporal variability of Src 1 is less significant than Src 15 and Src 28. This is likely because K-S tests are most sensitive around the median value of the independent variable. However, a visual examination of the lightcurve suggests that it is variable on timescales of about 1 hr or more, with flux variations of about $`50\%`$. We attempted Fourier analysis techniques to detect any potential periods on minute-hour timescales for Src 1, 15 and 28. No statistically meaningful periodicity is found in any of the sources. ### 4.3 X-Ray Hardness Ratios As the count rates of the resolved point sources are typically quite low directly analyzing the spectrum of the sources by fitting models is not practical. Instead we study the hardness ratios. Following Sarazin et al. (2000) and other authors, we measured the background-subtracted counts for each of the 47 resolved sources in three energy bands: soft (S), 0.3–1.0 keV; medium (M), 1.0–2.0 keV; and hard (H), 2.0–10.0 keV and calculated the hardness ratios R21 and R31 using the definitions R21 = (M $``$ S)/(M + S) and R31 = (H $``$ S)/(H + S), respectively. We list the calculated hardness ratios and their 1$`\sigma `$ errors, which are very large in general, in columns 8 and 9 of Table 1, and plot R31 versus R21 for all the sources in Figure 4. Only the typical 1$`\sigma `$ error bars are illustrated in order not to complicate the figure. In the same figure, we also show the predicted hardness ratios for an absorbed power-law model with different column densities $`N_\mathrm{H}`$ and photon indices $`\mathrm{\Gamma }`$. It can be seen in that most of the sources are located in a diagonal band from ($`0.217`$, $`0.600`$) to (+0.757, +0.837). We note that these colors are similar to those found in the X-ray bright elliptical galaxy NGC 720 (Jeltema et al. 2003) and NGC 4649 (Randall et al. Irwin 2004). Three sources, i.e., Src 17 (+0.678, +0.372), Src 31 (+0.757, +0.837), and Src 45 (+0.702, +0.436), appear to have been absorbed by a column density larger than the Galactic value. Since they are all located far away from the center of the galaxy ($`d>30^{\prime \prime }`$), they are probably unrelated background AGNs. However, we find that the hardness ratios of Src 17 and Src 45 cannot be reproduced by using a simple absorbed power-law spectral model, which is not in agreement with the spectrum of the hard X-ray cosmic background (Mushotzky et al. 2000). The brightest source in the field (Src 1) is located at the nucleus of the galaxy and has the X-ray color (R21, R31) = ($`0.086`$, $`0.504`$). Unlike in NGC 720 (Jeltema et al. 2003) and NGC 4649 (Randall et al. 2004), we found no supersoft sources (SSSs) in NGC 4552. Within the 4 $`R_\mathrm{e}`$ region the total X-ray emission (diffuse emission plus all the resolved point sources) of the galaxy has a hardness ratio of (R21, R31) = ($`0.453\pm 0.011`$, $`0.779\pm 0.015`$), as compared to the (R21, R31) = ($`+0.036\pm 0.031`$, $`0.221\pm 0.033`$) hardness ratio for the resolved sources and (R21, R31) = ($`0.565\pm 0.012`$, $`0.896\pm 0.018`$) for the unresolved diffuse emission from the gas plus unresolved point sources (see also §4.4). ### 4.4 Spectral Properties Within the 4 $`R_\mathrm{e}`$ region, there are 12 X-ray point sources (Src 1, 2, 6, 8, 11, 15, 24, 28, 37, 41, 43 and 46) each having more than 100 counts. We have extracted and studied the individual spectra of these bright sources. Among them, Src 1 is the brightest and coincides with the mini-AGN identified in the UV/optical bands (Cappellari et al. 1999). We will present detailed analysis of its spectrum separately in §4.6. We divide the rest of sources with less than 100 counts into four groups (Table 2). Sources in group A are the hardest ones that appeared at the top of Figure 4. Sources in group B are those for which the Galactic column density is required when an absorbed power-law model is applied to describe their hardness ratios, while for sources in group C a larger column density is needed. The rest of the sources are in group D. We have extracted and studied the cumulative spectra of each of the four groups, as well as that of all the resolved sources (excluding Src 1). In the model fittings, the background spectra were extracted from the local background fields adjacent to the regions where the source spectra were extracted. The fitting results are summarized in Table 2. As can be seen, excluding the central source, the cumulative spectrum of all resolved sources within the 4 $`R_\mathrm{e}`$ region can be fitted best by an absorbed power-law model with a photon index of $`\mathrm{\Gamma }=1.56\pm 0.07`$ when the column density is fixed to the Galactic value. The deduced intrinsic luminosity of all the resolved sources in 0.3–10 keV is $`1.83\pm 0.09\times 10^{40}`$ ergs s<sup>-1</sup>, assuming that the sources are all at the distance of NGC 4552. We also attempted to separate the emission of the unresolved point sources from the hot diffuse gas. In order to do this we extracted the spectrum of the total diffuse emission by excluding all resolved point sources. The background for this spectrum was extracted on the S3 chip as far away as possible from the galaxy. To fit the spectrum, we used a thermal component (apec model) to represent the contribution of the hot plasma, and a non-thermal component (power-law model) to represent the emission from the unresolved point sources, under the assumption that the hard spectral component seen in the diffuse emission is mainly due to the LMXBs. Both spectral components are subjected to a common absorption $`N_\mathrm{H}`$. Because we find that when $`N_\mathrm{H}`$ is free the obtained value is only slightly lower than, but still consistent with the Galactic value, we simply fix $`N_\mathrm{H}`$ to the Galactic value. We fixed the power-law photon index of the hard component to the value of 1.56 measured for the cumulative spectrum of all resolved sources; allowing the photon index to vary did not improve the fit. The best-fit gas temperature is found at $`kT=0.51\pm 0.02`$ keV, and the calculated luminosity of the total diffuse emission in 0.3–10 keV is $`L_{\mathrm{dif}}=4.32\pm 0.22\times 10^{40}`$ ergs s<sup>-1</sup>, of which about 31%, or $`1.34\pm 0.07\times 10^{40}`$ ergs s<sup>-1</sup>, can be ascribed to the unresolved point sources. If calculated in count flux, the contribution of the unresolved sources is about 20%. Thus, by adding the emissions of both resolved and unresolved point sources together, we estimate that in 0.3–10 keV the LMXBs contribute about 48% of the galaxy’s total luminosity, or about $`29\%`$ of the total count flux. Next we fit the cumulative spectra of the four groups of resolved point sources with an absorbed power-law model and list the results in Table 2. We find that for group A the obtained photon index is $`0.74\pm 0.27`$ if the column density is fixed to the Galactic value. The goodness of fit is marginally acceptable at $`\chi _r^2=15.7/13`$. When we allow the column density to vary, the fit can be improved significantly ($`\chi _r^2=7.2/12`$) in terms of the F-test which yields a best-fit photon index of $`1.83_{0.55}^{+0.85}`$, and a large absorption of $`N_\mathrm{H}=6.53_{2.10}^{+5.40}\times 10^{21}`$ cm<sup>-2</sup>. The spectra of group B and C can be well fit by the absorbed power-law model with $`\mathrm{\Gamma }=1.49_{0.17}^{+0.18}`$ and $`2.19_{0.32}^{+0.56}`$, respectively, when $`N_\mathrm{H}`$ is fixed to the Galactic value for group B and allowed to vary for group C. The spectrum of group D appears to be complicated and cannot be easily explained by using an one-component spectral model. ### 4.5 X-Ray Luminosity Function Based on the best-fit power-law model for the cumulative spectrum of all resolved sources within the 4 $`R_\mathrm{e}`$ region (excluding the central source; $`\mathrm{\Gamma }=1.56`$), we convert the observed count rates of the 46 off-center sources into unabsorbed 0.3–10 keV luminosities, assuming that the sources are all at the distance of NGC 4552. The conversion factor is $`5.18\times 10^{36}`$ ergs cts<sup>-1</sup>, and the resulting luminosities range from $`7\times 10^{37}`$ to $`1.5\times 10^{39}`$ ergs s<sup>-1</sup>. With these results we construct the XLF, in which the central source is not included. First, we fit the observed cumulative XLF that is not corrected for the effect of incompleteness at the faint end of the luminosity function. We use the maximum-likelihood method and determine the 90% confidence errors by performing Monte-Carlo simulations. Based on deep Chandra observations of blank fields (Mushotzky et al. 2000), we estimate that within the 4 $`R_\mathrm{e}`$ region of NGC 4552, $`{}_{}{}^{<}4`$ of the 47 detected point sources are expected to be unrelated background/foreground X-ray sources. In the fittings, we perform Monte-Carlo simulations to simulate the contribution of these background/foreground sources accordingly. We adopt either a single or a broken power-law model to fit the XLF. The single power-law profile is expressed as $$N(L_{38})=N_0L_{38}^\alpha ,$$ (1) where $`L_{38}`$ is the 0.3–10 keV luminosity in units of $`10^{38}`$ ergs s<sup>-1</sup>. The broken power-law profile can be expressed in the form of Eq.(1) at the high and low luminosities with different slopes $`\alpha _\mathrm{h}`$ and $`\alpha _\mathrm{l}`$, respectively. In the single power-law fitting, we obtain a slope $`\alpha =0.95_{0.2}^{+0.11}`$. However, the fit is poor with a probability of $`P=35\%`$ that the data and model are drawn from the same distribution. It overestimates the data at high luminosities and underestimates it at low luminosities. On the other hand, the broken power-law model can significantly improve the fit ($`P=74\%`$). The obtained break luminosity is at $`L_\mathrm{b}=3.7_{1.9}^{+1.4}\times 10^{38}`$ ergs s<sup>-1</sup>, and the two slopes are $`\alpha _\mathrm{h}=2.00_{0.71}^{+1.68}`$ and $`\alpha _\mathrm{l}=0.73_{0.41}^{+0.19}`$, respectively. Next we examine the effect of incompleteness at the faint end of the observed XLF by adopting a method similar to that outlined in Kim and Fabbiano (2004). With the use of the MARX package (Wise et al. 2003) we run a series of Monte-Carlo simulations to create fake X-ray point sources in NGC 4552 with given luminosities that cover the observed XLF’s range, and then determine how many of them can be detected with the same technique described in §4.1. We assume that the radial distribution of the fake sources at a given luminosity follows the standard r<sup>1/4</sup> law (de Vaucouleurs, 1948). With the obtained pick-out ratios at each given luminosity, we are able to correct both the observed XLF and the background. The results of the fittings for the corrected XLF are shown in Figure 5. The single power-law model obviously gives a better fit to the corrected XLF than to the uncorrected one with a steeper slope $`\alpha =1.18_{0.16}^{+0.13}`$ and a probability of 48% that the data and model are drawn from the same distribution. The broken power-law model gives the best fit ($`P=82\%`$), with $`L_\mathrm{b}=4.4_{1.4}^{+2.0}\times 10^{38}`$ ergs s<sup>-1</sup>, $`\alpha _\mathrm{h}=2.28_{0.53}^{+1.72}`$, and $`\alpha _\mathrm{l}=1.08_{0.33}^{+0.15}`$. We notice that at the 90% confidence level these best-fit parameters agree with those for the uncorrected XLF. For both the corrected and uncorrected XLFs, $`L_\mathrm{b}`$ and $`\alpha _\mathrm{h}`$ are relatively poorly determined due to the small number statistics. We also have attempted to fit the corrected XLF with a cutoff power-law model. Whether or not the incompleteness correction is made, we find that the model gives poor fit to the data and actually can be rejected at the 90% confidence level. These results agree with those in Kim and Fabbiano (2004), where NGC 4552 is not included in the sample. ### 4.6 Central Source Src 1 is the brightest X-ray source in the field. It has been speculated that it might be a mini-AGN based on observations in the optical/UV bands (Renzini et al. 1995; Cappellari et al. 1999). As previously noted in §3 and §4.2, our studies reveal that it is variable on the 1 hr and larger timescales, and its X-ray position agrees very well with the IR, optical, UV and radio centers of the galaxy. We have extracted the Chandra ACIS spectrum of this source and applied an absorbed power-law model to it. We find that the fit (Table 2) is acceptable with a best-fit photon index $`\mathrm{\Gamma }=2.11_{0.18}^{+0.19}`$. This confirms the result of Colbert and Mushotzky (1999), who obtained $`\mathrm{\Gamma }=2.24_{0.18}^{+0.20}`$ with ASCA data. Using the best-fit parameters, we estimate that in 0.3–10 keV the central source has an unabsorbed luminosity of $`3.99\pm 0.44\times 10^{39}`$ ergs s<sup>-1</sup>. Based on these multi-band parameters we conclude that the central source is most likely a low-luminosity AGN rather than a clump of LMXBs. ### 4.7 Off-Center Sources with $`L_\mathrm{X}>10^{39}`$ erg s<sup>-1</sup> We detected three off-center X-ray sources (Src 15, 41 and 43) that have luminosities larger than $`10^{39}`$ erg s<sup>-1</sup>. One of them (Src 15) is located at about $`26^{\prime \prime }`$ from the center of the galaxy, and has temporal variability on $`{}_{}{}^{>}1.5`$ hr timescales. The other two are at about $`90^{\prime \prime }`$ from the center with no evidence for variability during the observation. Src 41 is located in the joint Chandra-HST field and is found to be associated with a GC (§4.8). We extracted the ACIS spectrum for each of the three sources and fitted them with both an absorbed power-law model (PL) and an absorbed multicolor disk blackbody model (DBB). For Src 41, when the absorption is fixed to the Galactic value, the DBB model gives an acceptable fit ($`\chi _r^2=11.5/10`$) with a disk surface temperature at the inner radius of $`kT=0.75_{0.12}^{+0.14}`$ keV. The PL model gives a worse fit with a somewhat steep slope of $`\mathrm{\Gamma }=2.11_{0.19}^{+0.21}`$. If the absorption is left free, the goodness of the PL fit can be significantly increased, but the resulting absorption column density is about one order of magnitude larger than the Galactic value. For Src 15, both the PL model and DBB model give an acceptable fit. The best-fit power-law photon index and inner disk surface temperature are $`\mathrm{\Gamma }=1.25_{0.29}^{+0.32}`$. and $`kT=1.86_{0.62}^{+1.38}`$ keV, respectively, when the Galactic absorption value is used. For Src 43, when the absorption is fixed to the Galactic value, the DBB model gives a relatively poor fit ($`\chi _r^2=31.8/22`$) with an inner disk surface temperature of $`kT=1.17_{0.24}^{+0.32}`$ keV, while the PL model gives a slightly better fit ($`\chi _r^2=27.1/22`$) with a photon index of $`\mathrm{\Gamma }=1.62_{0.21}^{+0.17}`$. As both models have residuals at about 1 keV we allowed the absorption to be a free parameter and also introduced an additional soft spectral component in the fits. However, neither approach improved the fits. ### 4.8 ASSOCIATION WITH GLOBULAR CLUSTERS Kundu and Whitmore (2001) studied the globular cluster system of NGC 4552 in their analysis of 28 elliptical galaxies using the high spatial resolution images of the Wide Field Planetary Camera 2 (WFPC2) on board the HST. For NGC 4552 the field of view of the selected HST pointing lies entirely in the 4$`R_\mathrm{e}`$ region and covers an area of approximately 5 arcmin<sup>2</sup> (42% of the 4 $`R_\mathrm{e}`$ region), which includes the nucleus and the inner parts of the galaxy as well as a region out to about $`2^{}`$ southeast to the center (Fig.2). Twenty five of the off-center X-ray sources and 210 GC candidates with colors between $`0.5<VI<1.5`$ are detected in this field jointly covered by Chandra and HST. Hereafter, we restrict our analysis and discussions to these 25 X-ray sources and 210 GCs unless mentioned otherwise. There was a systematic offset of about $`2^{\prime \prime }`$ between the HST and Chandra positions. After correcting this offset by using the method outlined in Maccarone et al. (2003), we find that 10 X-ray sources match GCs within $`0.5^{\prime \prime }`$. Two other X-ray sources are within $`0.7^{\prime \prime }`$ and $`1.0^{\prime \prime }`$ of the nearest GCs, respectively, but they are not listed in the LMXB-GC matches in this work. Based on the spatial distributions of the X-ray sources and GCs, we estimate that $`{}_{}{}^{<}0.4`$ fake LMXB-GC matches should occur at random in the selected Chandra-HST field. Thus, we consider all of the 10 LMXB-GC matches to be real. The 10 LMXBs in GCs have count rates ranging from $`3.38\times 10^4`$ to $`5.74\times 10^3`$ cts s<sup>-1</sup>. They show no particular pattern of concentration in spatial distribution (Fig. 2). The one located nearest to the center of the galaxy is at a distance of $`15^{\prime \prime }`$. One of them (Src 41) is brighter than $`10^{39}`$ erg s<sup>-1</sup> (§4.7). #### 4.8.1 GCs Hosting Bright LMXBs and Non-LMXB GCs In Figure 6a and b, we compare the cumulative optical luminosity distributions of GCs that host bright LMXBs (LMXB GCs) and those that do not (non-LMXB GCs) as a function of their $`V`$-band and $`I`$-band magnitudes (Kundu & Whitmore 2001), respectively. As can clearly be seen, the detected X-ray sources tend to be associated preferentially with the optically bright GCs. In the $`V`$ band, the median value of the LMXB GC distribution is 22.0 mag, while the corresponding value of the non-LMXB GCs is 23.1 mag. Quite similar to this, in the $`I`$ band the two median values are 20.9 and 22.0 for LMXB GCs and non-LMXB GCs, respectively. Using the Mann-Whitney rank-sum tests (Mann & Whitney 1947), we calculate that the probability of the fraction distributions of LMXB GCs and non-LMXB GCs following the same distribution are only $`9\times 10^4`$ and $`3\times 10^4`$ in the $`V`$ and $`I`$ bands, respectively. Figure 6c shows the histogram of the $`VI`$ color distribution of all 210 GCs. The distribution is broad with a strong peak at $`VI0.95`$, and is more extended to the redder side. There is a rather weak structure at $`VI1.12`$, but the evidence for a bimodal distribution is not statistically significant (Kundu & Whitmore 2001). Such a distribution pattern implies the existence of a population of bluer, metal-poor GCs and a more dispersive population of redder, metal-rich GCs. This is consistent with the results found by other studies such as Neilsen and Tsvetanov (1999). We find that the $`VI`$ color distribution of LMXB GCs tends to have two corresponding concentrations. The number of LMXB GCs in the red population is $`1.9\pm 0.4`$ times as many LMXB GCs as in the blue population. To crosscheck this we plot the cumulative fraction distributions of LMXB GCs and non-LMXB GCs as a function of the $`VI`$ color (Fig. 6d). We see that the two distributions are different. The median colors of the LMXB GCs and non-LMXB GCs are $`VI=`$1.12 and 1.02, respectively. Using the Mann-Whitney rank-sum test, we find that the probability of the two distributions being the same is not significant (0.33). #### 4.8.2 GC-Associated LMXBs and Field LMXBs In Figure 7a we show both the cumulative number fraction of X-ray sources associated with GCs and that of the field sources located in the joint Chandra-HST field as a function of the intrinsic X-ray luminosity. We find that the median luminosity of the distribution is $`3.04\times 10^{38}`$ erg s<sup>-1</sup> for the GC-associated X-ray sources, and $`2.63\times 10^{38}`$ erg s<sup>-1</sup> for the field sources. The probability that the two distributions are the same is 0.50 with the Mann-Whitney rank-sum tests, which means that there is no significant difference between them. Thus X-ray sources in GCs and those not in GCs may have nearly the same mean X-ray luminosities. We study the X-ray spectral properties of GC-associated and field LMXBs, first by comparing their X-ray hardness ratios as has been defined in §4.3 and plotted in Figure 4, where the GC-associated sources are marked with open diamonds. However, it is hard to find any obvious difference in the distribution pattern directly on the X-ray color-color diagram. So we present the distributions of the source number as a function of the R21 and R31 colors for the GC-associated and field LMXBs, respectively (Fig. 7b and c). We find that the median values of the GC-associated sources are R21=$`+0.09`$ and R31$`=+0.06`$, which are slightly smaller than those of the field sources (R21=$`+0.12`$ and R31=$`+0.10`$). We thus speculated that the GC-associated LMXBs may be slightly softer than their counterparts in the field. In order to validate our speculation in a quantitative way, we fit the cumulative spectra of the GC-associated and field LMXBs by using an absorbed power-law model. The absorption column density is fixed to the Galactic value, because it cannot be well constrained if it is left free. When the 68% errors are quoted, the obtained photon indices are $`\mathrm{\Gamma }_{\mathrm{GC}}`$ = $`1.66\pm 0.08`$ for the GC-associated sources and $`\mathrm{\Gamma }_{\mathrm{Field}}`$ = $`1.42\pm 0.08`$ for the field sources, which supports that the GC-associated X-ray sources are softer. At the 90% confidence level, the model gives $`\mathrm{\Gamma }_{\mathrm{GC}}`$ = $`1.66\pm 0.12`$ and $`\mathrm{\Gamma }_{\mathrm{Field}}`$ = $`1.42_{0.13}^{+0.14}`$, which overlap only marginally. None of the GC-associated sources show significant temporal variability in X-rays. ## 5 DISCUSSION ### 5.1 The Break on the XLF The break on the XLF of point sources in early-type galaxies has been detected at $`35\times 10^{38}`$ ergs s<sup>-1</sup> in NGC 4697 (Sarazin et al. 2000), NGC 4472 (Kundu et al. 2002), NGC 1553 (Blanton et al. 2001), and a sample of 14 E/S0 galaxies (including NGC 4697; Kim and Fabbiano 2004). However, in NGC 720 (Jeltema et al. 2003) it is found at a much higher luminosity of $`L_\mathrm{b}=1.07_{0.1}^{+0.1}\times 10^{39}`$ ergs s<sup>-1</sup> for $`H_0`$ = 70 km s<sup>-1</sup> Mpc<sup>-1</sup>. Recent analyses of NGC 4365, NGC 4382 (Sivakoff et al. 2003), M49, M87, and NGC 4697 (Jord$`\stackrel{´}{a}`$n et al. 2004) suggest an upper cutoff at $`12\times 10^{39}`$ ergs s<sup>-1</sup>, rather than a break. In NGC 1600, Sivakoff et al. (2004) showed that a single power-law profile is sufficient to describe the observed XLF. In NGC 4552, we argue that a break at about $`4.4\times 10^{38}`$ ergs s<sup>-1</sup> is necessary to fit the XLF corrected for the sample incompleteness of faint sources (§4.5). For luminosities below $`5\times 10^{38}`$ erg s<sup>-1</sup>, the profile of the observed XLF is consistent with that expected for ultracompact binaries (Bildsten & Deloye 2004). Why are the conclusions on the XLF profile so dispersed from case to case? Is there an intrinsic universal break luminosity? Since the typical number of the detected X-ray sources is only 50-150 per galaxy we speculate that even if there is an universal break, the small number statistics will preclude us from measuring it. To investigate this possibility, we carry out direct Monte-Carlo simulations to create a series of fake XLFs whose profiles are intrinsically determined by a broken power-law profile with the typical parameters $`\alpha _\mathrm{l}=1.0`$, $`\alpha _\mathrm{h}=2.0`$ and $`L_\mathrm{b}=4.0\times 10^{38}`$ ergs s<sup>-1</sup>. For each simulated XLF, we consider 100 sources. By fitting the resulting XLFs with the broken power-law model, we find that at the 90% confidence level the obtained $`L_\mathrm{b}`$ ranges from $`2.8\times 10^{38}`$ ergs s<sup>-1</sup> to $`1.2\times 10^{39}`$ ergs s<sup>-1</sup>. Here the lower limit of $`L_\mathrm{b}`$ agrees well with the break luminosity found in galaxies such as NGC 4697, and the upper limit is consistent with the positions of the cutoffs found in NGC 720, NGC 4365, NGC 4382 and others. Such a large dispersion makes determination of the universal break, if any, observationally impossible. On the theoretical aspect, an universal break on the XLF is not necessary to signify the transition from neutron stars to black holes. The ultra compact binaries with He or C/O donors should have doubled Eddington luminosity comparing to their counterparts with H donors, and such systems are expected to exist in the dense GC environment in early-type galaxies (Bildsten & Deloye 2004). Also, LMXBs may emit above the Eddington limit if the emission is not isotropic. ### 5.2 Low Mass X-ray Binaries and Their Associations with Globular Clusters By analyzing the Chandra and HST data, we find that the fraction of the GCs hosting bright LMXBs in NGC 4552 (4.8%) is similar to those found in NGC 1399 (3.8%; Angelini et al. 2001), NGC 4472 (4%; Kundu et al. 2002), NGC 1553 (2.9%), NGC 4365 (5.5%), NGC 4649 (4.9%) and NGC 4697 (2.7%; Sarazin et al. 2003). Moreover, the fraction of LMXBs associated with GCs is 40% in NGC 4552, which is similar to those of the X-ray bright ellipticals NGC 4472 (40%; Kundu et al. 2002) and NGC 4649 (47%), and the X-ray faint ellipticals NGC 4365 (49%) and NGC 4697 (44%; Sarazin et al. 2003), which is consistent with the argument of White et al. (2002) that at present GCs are the dominant sites for LMXB formation in early-type galaxies. The fraction of GC-associated LMXBs in NGC 4552 is higher than those of the S0 galaxies NGC 1553 (18%; Sarazin et al. 2003) and NGC 1332 (30%; Humphrey et al. 2004), but smaller than that of the cD galaxy NGC 1399 (70%; Angelini et al. 2001). Considering that in typical spiral galaxies such as our Galaxy and M31 (e.g., Supper et al. 1997) the fraction is only about 10%, these observations are generally in agreement with previous suggestions (e.g. Sarazin et al. 2003) that the fraction of X-ray sources residing in GCs may increase along the Hubble sequence from spiral bulges to S0, E, and then cD galaxies. However, it is important to note that the LMXB-GC connection has only been studied in small central regions observed by the HST in each of the early type galaxies mentioned here. This apparent variation with Hubble type may be amplified by any spatial variation in the rate of GC-LMXB associations. We note that considering only the LMXBs in the inner few kpc of the Milky Way to the comparable HST based analyses of the more distant early type galaxies would suggest similar GC-LMXB association rates in early and late type galaxies. We calculate that in NGC 4552 the probability of having bright X-ray sources in GCs is $`1.47\times 10^7`$ LMXB per $`L_{,I}`$, which agrees very well with the values obtained by Kundu et al. (2003) and Sarazin et al. (2003). Bildsten & Deloye (2004) argue that this can be explained with ultracompact binaries that have a birthrate of one new mass transferring binary every $`2\times 10^6`$ yr per $`10^7M_{}`$ of GCs. We find that GCs hosting bright LMXBs are typically 1–2 magnitudes brighter than those with no detected LMXBs in the $`V`$ and $`I`$ bands. In fact at the level of significance of $`10^3`$ we reject the hypothesis that the luminosities of LMXB GCs and non-LMXB GCs are drawn from the same distribution. We also find that in the red, metal-rich GCs there are about $`1.9\pm 0.4`$ times as many LMXBs as there are in the blue, metal-poor ones. These results are consistent with, or quite similar to those found in NGC 4472 (Kundu et al. 2002), M87 (Jord$`\stackrel{´}{a}`$n et al. 2004) and other S0/E galaxies (Sarazin et al. 2003), indicating that the high GC formation efficiency is largely attributed to the metallicity, rather than the age of the old stellar systems (Kundu et al. 2003). In NGC 4472 and M87, the number of redder GCs that host a LMXB is about 3 times more than their bluer counterparts. The relatively lower overabundance of LMXB ($`2`$) in redder GCs in NGC 4552 may be partly due to the fact that the redder GC population is less prominent in this galaxy (Kundu & Whitmore 2001). A comparison of the X-ray properties of the GC LMXBs and field populations of LMXBs may also help us understand the origin of the field LMXBs. In NGC 4552, we find no obvious difference between the luminosity distributions of GC and field LMXBs. By examining both the hardness ratios and the cumulative X-ray spectra of the GC and field sources, we find that the LMXBs in GCs are softer than the field LMXBs at the 68% confidence level. Neither the GC LMXB sample nor field LMXBs has particularly bright soft or hard sources that may dominate the counts in the combined spectrum, which can bias the analysis significantly. Thus, the spectral difference appears to be physically real. This may imply that in NGC 4552 the mean metallicity of the GC-associated LMXBs is higher than that of the field LMXBs (Maccarone et al. 2004), as opposed to NGC 4472, where the LMXBs in GCs tend to be slightly harder (Maccarone et al. 2003). ### 5.3 Brightest Off-Center X-Ray Sources Typically, 1–4 X-ray sources that have a 0.3–10 keV luminosity larger than $`10^{39}`$ erg s<sup>-1</sup> are detected in each early-type galaxy studied to date (e.g., Sarazin et al. 2001; Blanton et al. 2001; Angelini et al. 2001; Kim & Fabbiano 2003; Humphrey & Buote 2004). In a few cases, such as NGC 720 (Jeltema et al. 2003) and NGC 1600 (Sivakoff et al. 2004), an even larger number has been reported. The nature of these very bright sources is not clear. Some of these studies suggest that they are ultra-luminous X-ray sources (ULXs) that host an intermediate mass black hole (IMBH). However, in a recent study of nearby galaxies Irwin et al. (2004) showed that in early type galaxies the sources brighter than $`2\times 10^{39}`$ erg s<sup>-1</sup> are most probably unassociated with the galaxy, while the sources with luminosities of $`12\times 10^{39}`$ erg s<sup>-1</sup> can be explained by accretion onto 10-20 $`M_{}`$ stellar mass black holes (SMBHs). Also in NGC 720, when a more conservative distance to the galaxy is adopted the number of the very bright sources associated with the galaxy is not statistically significant. The three brightest off-center X-ray sources detected within the 4 $`R_\mathrm{e}`$ region of NGC 4552 have isotropic 0.3–10 keV luminosities of $`1.18\times 10^{39}`$ ergs s<sup>-1</sup> (Src 15), $`1.15\times 10^{39}`$ ergs s<sup>-1</sup> (Src 41) and $`1.54\times 10^{39}`$ ergs s<sup>-1</sup> (Src 43), respectively. One of them, Src 41, is in the joint Chandra-HST field and is found to be associated with a globular cluster. Similar very bright X-ray source-GC matches also have been seen in other early-type galaxies by Angelini et al. (2001; NGC 1399) and Jeltema et al. (2003; NGC 720). We speculate that these very bright GC-associated sources may be powered by the accretion onto a black hole that is at the center of the host GC. Actually the X-ray spectrum of Src 41 in NGC 4552 can be better fitted with a single multiple blackbody disk model than a power-law model (assuming Galactic absorption). The best-fit inner disk temperature ($`0.75_{0.12}^{+0.14}`$ keV) is higher than that of NGC X-1 ($`kT_{\mathrm{in}}0.1`$ keV; Colbert & Mushotzky 1999), M81 X-9 ($`kT_{\mathrm{in}}0.2`$ keV; Miller et al. 2004) and others, inferring that the accreting source in Src 41 should have a relatively low mass, which is estimated to be about 23 $`M_{}`$ for a Schwarzschild black hole, or 15–135 $`M_{}`$ for a Kerr black hole, depending on the black hole spin and the sense of the disk rotation. This ambiguous result makes it difficult to distinguish between the IMBH and SMBH natures for Src 41. Although it is also possible that Src 41 is a background AGN, the probability for such an AGN-GC match is very low. For Src 43, when the absorption is fixed to the Galactic value, the multiple blackbody disk model gives an unacceptable fit to its X-ray spectrum. A better and marginally acceptable fit can be obtained with the power-law model with a photon index of $`1.62_{0.21}^{+0.17}`$. So this source is likely to be a neutron star binary system with beamed emissions or a SMBH binary system. For Src 15, both the power-law model and multiple blackbody model give acceptable fits. The best-fit inner disk surface temperature ($`1.86_{0.62}^{+1.38}`$ keV) is consistent with that of a neutron star or a SMBH system. Still, we cannot exclude the possibility that these two source are background AGNs. ## 6 SUMMARY By analyzing Chandra ACIS data we have detected 47 X-ray point sources within the inner 4 $`R_\mathrm{e}`$ region of the early-type galaxy NGC 4552. Most of the sources are inferred to be LMXBs. The position of the brightest point source is consistent with that of the IR, optical, UV and radio centers of the galaxy. In the X-ray band, the central source shows a relatively steep power-law spectrum and temporal variability on $`{}_{}{}^{}{}_{}{}^{>}`$ 1 hr timescales. These results confirm the early speculation that a low-luminosity AGN resides in the center of this galaxy (Renzini et al. 1995; Cappellari et al. 1999). The derived 0.3–10 keV luminosities of the 46 off-center sources range from $`7\times 10^{37}`$ to $`1.5\times 10^{39}`$ erg s<sup>-1</sup>. Three sources have isotropic 0.3–10 keV luminosities larger than $`10^{39}`$ erg s<sup>-1</sup>. One of them (Src 41) is in the joint Chandra-HST field and is associated with a globular cluster, and another (Src 15) shows temporal variations on $`{}_{}{}^{>}1.5`$ hr timescales. By studying their ACIS spectra we find that Src 41 may be a black hole system with a mass of 15–135 $`M_{}`$, while the other two sources should have lower masses if they are associated with the galaxy and not background AGN. We find that after correcting for the incompleteness at the low luminosity end the observed cumulative XLF can be best fit by a broken power-law model with a break at $`L_\mathrm{b}=4.4_{1.4}^{+2.0}\times 10^{38}`$ ergs s<sup>-1</sup>, while the single power-law model and the cutoff power-law model give worse, and unacceptable, fits. The position of the break is consistent with that found by Kim and Fabbiano (2004) in a sample of 14 E/S0 galaxies. By performing Monte-Carlo simulations we argue that even if there is an universal break, it is not a reliable distance indicator due to small number statistics. In an area jointly covered by both the Chandra ACIS and HST WFPC2, we detected 25 off-center X-ray point sources and 210 GCs, including 10 LMXB-GC matches. We find that the fraction of the GCs hosting bright LMXBs (4.8%) and the fraction of LMXBs associated with GCs (40%) are both in good agreement with those in other early-type galaxies (e.g., Kundu et al. 2002; Sarazin et al. 2003). As in NGC 4472 (Kundu et al. 2002) and M87 (Jord$`\stackrel{´}{a}`$n et al. 2004) and other early type galaxies, in NGC 4552 the GCs hosting bright LMXBs are typically 1–2 magnitudes brighter than the GCs with no detected LMXBs in the $`V`$ and $`I`$ bands. Moreover, there are about $`1.9\pm 0.4`$ times as many LMXBs in the red, metal-rich GCs as there are in the blue, metal-poor ones. This supports the idea that the high GC formation efficiency is largely attributed to the metallicity in old stellar systems (Kundu et al. 2003). We find no obvious difference between the X-ray luminosity distributions of GC LMXBs and field LMXBs. The cumulative spectrum of the LMXBs in GCs tend to be softer than that of the field LMXBs, which differs from result of Maccarone et al. (2003) who showed that in NGC 4472 the LMXBs in GCs tend to be slightly harder. This may indicate that in this galaxy the mean metallicity of the GC-associated LMXBs is higher than that of the field LMXBs (Maccarone et al. 2004). This work was supported by the National Science Foundation of China (Grant No. 10273009 and 10233040), and by the Ministry of Science and Technology of China, under Grant No. NKBRSF G19990754. AK thanks NASA for support via LTSA grant NAG5-12975.
warning/0506/quant-ph0506154.html	ar5iv	text	# No-Flipping as a consequence of No-Signalling and Non-increase of Entanglement under LOCC ## Abstract Non existence of Universal Flipper for arbitrary quantum states is a fundamental constraint on the allowed operations performed on physical systems. The largest set of qubits that can be flipped by a single machine is a great circle of the Bloch-sphere. In this paper, we show the impossibility of universal exact-flipping operation, first by using the fact that no faster than light communication is possible and then by using the principle of “non-increase of entanglement under LOCC”. Interestingly, in both the cases, there is no violation of the two principles if and only if the set of states to be flipped, form a great circle. PACS number(s): 03.67.Mn, 03.67.Hk Keywords: Flipping, Entanglement, No-signalling. <sup>1</sup>Department of Applied Mathematics, University of Calcutta, 92 A. P. C. Road, Kolkata- 700 009, India <sup>2</sup>Physics and Applied Mathematics Unit, Indian Statistical Institute, 203 B. T. Road, Kolkata- 700 108, India The structure of the allowed operations performed on the quantum systems imposes some restrictions on the systems. Sometimes these restrictions play a crucial role to understand the basic features of the system and naturally our task is to find the fundamental nature of these restrictions in a simple way. It has been shown that an arbitrary state taken from a set of two known, non-orthogonal states can not be copied exactly in a deterministic way . Similarly, one can not delete copy of an unknown state by performing some linear, trace preserving joint operations on two copies of that state . Interestingly, several authors derived these no-cloning and no-deleting theorems by applying some fundamental principles of nature like impossibility of signalling , preservation of entanglement for closed systems under local operations or rather increase of entanglement by LOCC . Another interesting feature of quantum system, is the non-existence of universal flipping machine for arbitrary input qubit states, i.e., there exists no universal flipper which can operate on any unknown qubit state $`\|\psi `$ resulting the orthogonal state $`\|\psi ^{}`$ . This no-flipping theorem has a stark dissimilarity with others, as unlike no-cloning and no-deleting, two non-orthogonal states can always be flipped. Actually the largest set of states (of qubit system) which can be flipped exactly, by a single unitary operator is the set of states lying on a great circle of the Bloch sphere . In this paper our aim is to establish the no-flipping theorem by applying the following established principles of nature: 1. Impossibility of superluminal signalling - It is impossible to communicate any message between some spatially separated parties with a speed greater than the speed of light. 2. The thermodynamical law of Entanglement - Amount of Entanglement shared between some spatially separated parties can not be increased by LOCC, i.e.,* by performing local operations on the subsystems and classical communications between them.* In other words our aim is to show, if exact flipping of even the minimal number (i.e., three) of states, not taken from one great circle, is possible, then one can send instantaneous signal as well as increase entanglement between two distant parties by local operations. For this purpose we consider three arbitrary states not lying in one great circle in their simplest form as; > $$\begin{array}{ccc}\|0,\hfill & & \\ \|\psi =a\|0+b\|1,\hfill & & \\ \|\varphi =c\|0+de^{i\theta }\|1,\hfill & & \end{array}$$ > (1) where a, b, c, d are real numbers satisfying the relation $`a^2+b^2=1=c^2+d^2`$ and $`0<\theta <\pi ,a>0,c>0`$, and the states $`\|0`$ , $`\|1`$ are orthogonal to each other. We assume that a machine exists which can flip at least these three states exactly. The most general flipping operation for these three states can be described as $$\begin{array}{ccc}& & \|0\|M\|1\|M_0\hfill \\ & & \|\psi \|Me^{i\mu }\|\overline{\psi }\|M_\psi \hfill \\ & & \|\varphi \|Me^{i\nu }\|\overline{\varphi }\|M_\varphi \hfill \end{array}$$ (2) where $`\mu `$ and $`\nu `$ are some arbitrary phases and $`\|M`$ is the initial machine state. The flipped states are orthogonal to the original states, i.e., $$\begin{array}{ccc}& & 0\|1=\psi \|\overline{\psi }=\varphi \|\overline{\varphi }=0,\hfill \end{array}$$ (3) where $`\|\overline{\psi }=b\|0a\|1`$ and $`\|\overline{\varphi }=de^{i\theta }\|0c\|1`$ in their usual notations. The above operation is not assumed to be unitary and the operation acts linearly on one side of an entangled state only when the density matrix has a mixture representation of the states in equation $`(2)`$. First we show that the exact flipping machine which we have considered, implies signalling. Consider two spatially separated parties, say, Alice and Bob who initially share an entangled state of the form, $$\begin{array}{ccc}\|\mathrm{\Psi }_{AB}^i\hfill & =& \frac{1}{\sqrt{3}}(\|0_A\|0_B+\|1_A\|\psi _B+\|2_A\|\varphi _B)\|M_B\hfill \end{array}$$ (4) where Alice holds a system associated with three dimensional Hilbert space, having a basis, $`\{\|0,\|1,\|2\}`$(say) and Bob’s system consists of a qubit (entangled with Alice’s system) and a flipping machine defined as in equation (2). Here one should note that the joint system of Alice and Bob has been chosen in such a manner that the marginal density matrix of Bob’s side admits a representation in terms of the three states $`\|0,\|\psi ,\|\varphi `$, on which the flipping machine has been defined. The reduced density matrix of Alice’s side is $$\begin{array}{ccc}\rho _A^i\hfill & =& \frac{1}{3}\{P[\|0]+P[\|1]+P[\|2]+a(\|01\|+\|10\|)\hfill \\ & & +c(\|02\|+\|20\|)+\varphi \|\psi \|12\|+\psi \|\varphi \|21\|\}\hfill \end{array}$$ (5) Now assume that Bob applies the flipping machine on his qubit. After the flipping operation the shared state between Alice and Bob takes the following form, $$\begin{array}{ccc}\|\mathrm{\Psi }_{AB}^f=\frac{1}{\sqrt{3}}\{\|0_A\|1M_0_B+e^{i\mu }\|1_A\|\overline{\psi }M_\psi _B+e^{i\nu }\|2_A\|\overline{\varphi }M_\varphi _B\}\hfill & & \end{array}$$ (6) The final density matrix of Alice’s side (expanded in the computational basis) is $$\begin{array}{ccc}\rho _A^f\hfill & =& \frac{1}{3}\{P[\|0]+P[\|1]+P[\|2]a(e^{i\mu }M_\psi \|M_0\|01\|\hfill \\ & & +e^{i\mu }M_0\|M_\psi \|10\|)c(e^{i\nu }M_\varphi \|M_0\|02\|+e^{i\nu }M_0\|M_\varphi \|20\|)\hfill \\ & & +\psi \|\varphi e^{i(\mu \nu )}M_\varphi \|M_\psi \|12\|+\varphi \|\psi e^{i(\nu \mu )}M_\psi \|M_\varphi \|21\|\}\hfill \end{array}$$ (7) As this is a trace preserving local operation performed entirely on Bob’s side and there is also no classical communication between them, so to prevent any violation of the principle of no-signalling, the reduced state on Alice’s side must remain unchanged. Equating the reduced density matrices on Alice’s side before and after the flipping operation on Bob’s side we get the following relations $$a=ae^{i\mu }M_0\|M_\psi =ae^{i\mu }M_\psi \|M_0$$ (8) $$c=ce^{i\nu }M_0\|M_\varphi =ce^{i\nu }M_\varphi \|M_0$$ (9) $$\varphi \|\psi =e^{i(\mu \nu )}\psi \|\varphi M_\varphi \|M_\psi $$ (10) $$\psi \|\varphi =e^{i(\nu \mu )}\varphi \|\psi M_\psi \|M_\varphi $$ (11) The above relations imply, $$\text{either}b=0\text{or}d=0\text{or}\mathrm{sin}\theta =0,$$ forcing the states to lie on a great circle. So, we observe that as long as the three states do not lie on a great circle, we have, $`\rho _A^i\rho _A^f`$, and interestingly whenever the equality holds, i.e., $`\rho _A^i=\rho _A^f`$, the three states actually lie on a great circle. This clearly shows that exact flipping of any three states not lying on a great circle is an impossibility. Now, we are going to show the impossibility of universal flipping using the principle of non-increase of entanglement by LOCC. Here we must be careful about the fact that in the earlier set-up (equation $`(4)`$), eigen values of $`\rho _A^i`$ and $`\rho _A^f`$ are equal, implying no change in entanglement even when states are not on the great circle. For this we choose another state, $`\|\mathrm{\Psi }_{AB}^i`$ of five qubits, shared between Alice and Bob, situated at distant locations, where the first qubit is with Alice, and remaining four($`B_1,B_2,B_3,B_4`$) are with Bob. $$\begin{array}{ccc}\|\mathrm{\Psi }_{AB}^i\hfill & =& \frac{1}{\sqrt{8}}\{(\|000+\|111)_{AB_1B_2}\|10_{B_3B_4}\hfill \\ & & (\|010+\|100+\|101)_{AB_1B_2}\|\overline{\psi }\psi _{B_3B_4}\hfill \\ & & (\|011+\|110+\|001)_{AB_1B_2}\|\overline{\varphi }\varphi _{B_3B_4}\}\|M_{B_M}\hfill \end{array}$$ (12) Let us assume that Bob has a flipping machine (defined earlier) with him which he applies on his last qubit($`B_4`$). After the flipping operation the joint state between them takes the form, $$\begin{array}{ccc}\|\mathrm{\Psi }_{AB}^f\hfill & =& \frac{1}{\sqrt{8}}\{(\|000+\|111)_{AB_1B_2}\|11M_0_{B_3B_4B_M}\hfill \\ & & e^{i\mu }(\|010+\|100+\|101)_{AB_1B_2}\|\overline{\psi }\overline{\psi }M_\psi _{B_3B_4B_M}\hfill \\ & & e^{i\nu }(\|011+\|110+\|001)_{AB_1B_2}\|\overline{\varphi }\overline{\varphi }M_\varphi _{B_3B_4B_M}\}.\hfill \end{array}$$ (13) In order to compare the amount of entanglement present in $`\|\mathrm{\Psi }_{AB}^i`$ with that present in $`\|\mathrm{\Psi }_{AB}^f`$, we compare the eigenvalues of the marginal density matrices on Alice’s side before and after the flipping operation. The initial state of Alice’s subsystem is; $$\begin{array}{ccc}\rho _A^i\hfill & =& \frac{1}{8}\{4(P[\|0]+P[\|1])+(a^2+c^2+2\|\psi \|\varphi \|^2)(\|01\|+\|10\|)\}.\hfill \end{array}$$ (14) The final state of Alice’s subsystem is; $$\begin{array}{ccc}\rho _A^f\hfill & =& \frac{1}{8}\{4(P[\|0]+P[\|1])+(2Za^2X^{}c^2Y)\|01\|\hfill \\ & & +(2Za^2Xc^2Y^{})\|10\|\}\hfill \end{array}$$ (15) where, $`Z=Re[e^{i(\mu \nu )}(\psi \|\varphi )^2M_\varphi \|M_\psi ],X=e^{i\mu }M_0\|M_\psi `$ and $`Y=e^{i\nu }M_0\|M_\varphi .`$ After constructing the eigenvalue equations for both $`\rho _A^i`$ and $`\rho _A^f`$ we find that the largest eigenvalues are respectively $$\lambda ^i=\frac{1}{2}+\frac{1}{8}(2\|\psi \|\varphi \|^2+a^2+c^2)$$ and $$\lambda ^f=\frac{1}{2}+\frac{1}{8}\|2Za^2X^{}c^2Y\|$$ It can easily be checked that $`\lambda ^f\lambda ^i`$( see appendix). This implies, $`E(\|\mathrm{\Psi }_{A:B}^f)E(\|\mathrm{\Psi }_{A:B}^i)`$ where $`E`$ stands for amount of entanglement. So this establishes the impossibility of universal flipping as here the only operation that has been allowed is local. A close look will reveal that the greatest eigenvalues are equal (i.e., no increase of entanglement), only when, the three states on which we have defined our flipping machine, lie on one great circle. The above construction seems to be a complicated one due to our linearity assumption. We have considered our flipping operation acts linearly on one side of an entangled state only when the reduced density matrix has a mixture representation of the states in equation $`(2)`$. However a more simpler proof is possible consisting only a three qubit system if we allow further that the operation is linear on superposition level. Now consider a three qubit state shared between Alice and Bob where Bob has a two qubit system ($`B_1`$ and $`B_2`$) as follows, $$\begin{array}{ccc}\|\mathrm{\Phi }^i_{AB}=\frac{1}{\sqrt{b^2+d^2}}\{\|0_A\frac{\|0_{B_1}\|\psi _{B_2}\|\psi _{B_1}\|0_{B_2}}{\sqrt{2}}+\|1_A\frac{\|0_{B_1}\|\varphi _{B_2}\|\varphi _{B_1}\|0_{B_2}}{\sqrt{2}}\}\hfill & & \\ =\frac{1}{\sqrt{b^2+d^2}}\{(b\|0_A+de^{i\theta }\|1_A)\frac{\|0_{B_1}\|1_{B_2}\|1_{B_1}\|0_{B_2}}{\sqrt{2}}\}\hfill & & \end{array}$$ (16) Clearly $`\|\mathrm{\Phi }^i_{AB}`$ is a separable (pure product) state in $`A:B`$ cut. Assume Bob has a flipping machine with him which he applies on his last qubit, i.e., on $`B_2`$. Then the joint state between them takes the form, $$\begin{array}{ccc}\|\mathrm{\Phi }^f_{AB}=\frac{1}{\sqrt{N}}\{e^{i\mu }\|00\|\overline{\psi }\|M_\psi +e^{i\nu }\|10\|\overline{\varphi }\|M_\varphi (\|0\|\psi +\|1\|\varphi )\|1\|M_0\}\hfill & & \end{array}$$ (17) where $`N=2+a^2Re\{e^{i\mu }M_\psi M_0\}+c^2Re\{e^{i\nu }M_\varphi M_0\}`$. It is easy to check that the state in general is an entangled state in $`A:B`$ cut, thus in this setting we also observe an increase of entanglement by LOCC. In the conclusive remarks, we want to mention that, as a constraint on quantum mechanical system, the ‘No-Flipping’ theorem is weaker than the ‘No-Cloning’ and ‘No-Deleting’ theorem, because exact flipping is possible for all states lying in one great circle of the Bloch sphere. Hence, showing the impossibility of universal flipping from the principle like No-signalling and Non increase of entanglement by local operations, is an interesting problem to deal with. We find the nonphysical nature of flipping operation on whole Bloch sphere in two different ways. In this context, recently no-flipping has been established by using the existence of incomparable states where one qutrit and two qubits are required . Appendix To check that, $`\lambda ^f\lambda ^i`$ for the second set-up, we have, $`\lambda ^i\lambda ^f`$ $`2\|\psi \|\varphi \|^2+a^2+c^2\|2Za^2X^{}c^2Y\|`$ $`a^4(1\|X\|^2)+c^4(1\|Y\|^2)+2a^2c^2[1Re(XY)]+4a^2[\|\psi \|\varphi \|^2+ZRe(X)]+4c^2[\|\psi \|\varphi \|^2+ZRe(Y)]+4[\|\psi \|\varphi \|^4Z^2]0`$ As, all terms on the L.H.S. of the last inequality are non-negative. Therefore, the above inequality for eigenvalues is satisfied. The equality holds only when, $`\|X\|=\|Y\|=1,Re(XY)=1,\|\psi \|\varphi \|^2+ZRe(X)=0,\|\psi \|\varphi \|^2+ZRe(Y)=0,Z^2=\|\psi \|\varphi \|^4.`$ Now, $`Z=Re[e^{i(\mu \nu )}(\psi \|\varphi )^2M_\varphi \|M_\psi ],X=e^{i\mu }M_0\|M_\psi `$ and $`Y=e^{i\nu }M_0\|M_\varphi ,`$ which implies machine states will differ by only some phases and the states $`\|0,\|\psi ,\|\varphi `$ will lie on a great circle. Hence we see that in this case, where $`\rho _A^i=\rho _A^f`$, the three states of equation (1) on which we have defined our flipping machine, actually lies on a great circle of the Bloch sphere. Acknowledgement. We thank Prof. N. Gisin for his valuable comments and suggestions. Also we acknowledge the referee for his/her valuable comments and suggestions. I.C. and S.K. also acknowledges CSIR, India for providing fellowship during this work.
warning/0506/hep-th0506136.html	ar5iv	text	# Plate with a hole obeys the averaged null energy condition ## I Introduction The standard Casimir calculation of the energy density between a pair of parallel plates (see for example Mostepanenko and Trunov (1997)) yields a negative energy density between the plates. While this result poses no problem in the calculation of the usual Casimir force, it presents a puzzle for general relativity. One can construct a spacetime with an arbitrary geometry $`R_{\lambda \nu }`$ simply by constructing the energy-momentum tensor to solve Einstein’s equations $$T_{\lambda \nu }=\frac{1}{8\pi G}\left(R_{\lambda \nu }\frac{1}{2}g_{\lambda \nu }R\right).$$ (1) The only way to prevent the appearance of exotic phenomena, such as closed timelike curves Hawking (1992), traversable wormholes Morris et al. (1988), or superluminal travel Olum (1998), is to place restrictions on the allowed energy-momentum tensors $`T_{\lambda \nu }`$. While these conditions are all obeyed in classical physics, the negative energy density of the quantum Casimir system violates most such conditions, including the weak energy condition (WEC) and the null energy condition (NEC), which require that $`T_{\lambda \nu }V^\lambda V^\nu 0`$ for timelike and null vectors respectively. A still weaker condition, which is still strong enough to rule out exotic phenomena (and to prove singularity theorems Penrose (1965); Galloway (1981); Roman (1986, 1988)), is that the null energy condition hold only when averaged over a complete geodesic (ANEC).<sup>1</sup><sup>1</sup>1For singularity theorems, the average is only over the future of the trapped surface. Geodesics parallel to the plates obey NEC, so any candidate for ANEC violation would need to pass through the plates themselves. Therefore one cannot test ANEC using the standard Casimir calculation in which one imposes ideal boundary conditions, since this calculation is not valid within each plate. One approach to resolve this question is to model the plate using a domain wall background Graham and Olum (2003); Olum and Graham (2003); in this case the effect of the domain wall modifies the calculation significantly, so that ANEC is obeyed. In a number of other examples in which one might expect to find that ANEC is violated, explicit calculation shows that it is obeyed Schwartz-Perlov and Olum (2003); Graham et al. (2004). Other calculations also show that energy condition violation is more difficult to achieve in realistic situations than idealized models would suggest Sopova and Ford (2002, 2005). ANEC is also known to be obeyed by free scalar Klinkhammer (1991) and electromagnetic Folacci (1992) fields in flat spacetime. Other works have found restrictions on energy condition violation in flat space Borde et al. (2002); Ford and Roman (1995, 1996). In this paper we consider an alternative modification to the Casimir problem that one might expect would allow the NEC violation between the plates to extend to ANEC violation: we imagine a plate with a small hole, through which the geodesic can pass without encountering the material of the plate. Our primary calculation is for the case of a single Dirichlet plate, which also leads to a negative energy density for minimal coupling. We consider both two and three spatial dimensions. We also give extensions of this result to Neumann boundaries and to the case of two plates in extreme limits. ## II Null Energy Condition Outside a Boundary Since our geodesic never passes through the material that actually imposes the boundary, we need only consider the quantum field $`\varphi `$ in empty space. The effect of the boundary will be to modify the normal mode expansion for $`\varphi `$. We will then integrate $`V^\lambda V^\nu T_{\lambda \nu }`$, where $`T_{\lambda \nu }`$ is the stress-energy tensor, over the null geodesic $`V^\lambda `$ perpendicular to the plate, passing through the center of the hole. The stress-energy tensor for a minimally-coupled scalar field is $$T_{\lambda \nu }=_\lambda \varphi _\nu \varphi \frac{1}{2}\eta _{\lambda \nu }\left[^\lambda \varphi _\lambda \varphi \right].$$ (2) For a null vector, $`\eta _{\lambda \nu }V^\lambda V^\nu =0`$, so we have $$T_{\lambda \nu }V^\lambda V^\nu =\left(V^\alpha _\alpha \varphi \right)^2.$$ (3) For a static system, $`T_{0i}=0`$ for $`i=1,2,3`$. If we further choose spatial coordinates in which $`T_{ij}`$ is diagonal and $`V=(1,𝐯)`$, then $$T_{\lambda \nu }V^\lambda V^\nu =\dot{\varphi }^2+\underset{i}{}\left(v_i_i\varphi \right)^2.$$ (4) Let the center of the hole lie at the origin and let the $`z`$ axis be the direction perpendicular to the plate, along which the geodesic lies. For that path, $$T_{\lambda \nu }V^\lambda V^\nu =\dot{\varphi }^2+\left(_z\varphi \right)^2.$$ (5) ## III Babinet’s Principle It will greatly simplify the scattering theory techniques we would like to use in our Casimir calculation Bordag and Lindig (1996); Saharian (2001); Graham et al. (2002) to be able to consider a boundary condition in a local region. Therefore we apply a Babinet’s principle argument to reexpress the result for a plate with a hole in terms of the results when the boundary condition is applied to an entire plate and when a complementary boundary condition is applied to a disk. The former is well-known, while the latter can be computed using scattering theory in elliptical or spheroidal coordinates. We start in empty space, and write the field there in terms of normal modes that are even or odd in the coordinate across the boundary. Next we consider a perfectly reflecting Dirichlet boundary with no holes. The free-space odd modes obey the boundary conditions, but the even modes do not. Instead, we have new even modes, which are just the odd mode on the right and minus the odd mode on the left, as shown in Fig. 1. If we let $`E`$ denote a sum over the free-space even modes and $`O`$ the same sum over the odd modes, then in free space we have $`E+O`$, whereas with the barrier we have $`O+O`$. Therefore, the renormalized energy, the difference between the energy with the barrier and the energy in free space, is $`OE`$. If we have Neumann conditions instead, the situation is precisely reversed and the energy is $`EO`$. Now we consider the Dirichlet plate with holes of arbitrary shapes. Once again the odd modes are unaffected, and we have new even modes that vanish on the barrier but are continuous in the holes, as shown in Fig. 2. Since they are even, they satisfy Neumann conditions in the hole. Let us call the contribution of those modes $`A`$. The energy with the perforated barrier is thus $`A+O`$, so the renormalized energy is $`AE`$. Finally, suppose that there are Neumann patches where the holes were. The even modes are unaffected, but there are new odd modes. In order to be odd and continuous they must satisfy Dirichlet conditions on the plane outside the patches. Thus, except for a change of sign on one side, these are the exact same modes of the previous paragraph, as shown in Fig. 2. Therefore the total energy is $`A+E`$ and the renormalized energy is $`AO`$. Thus we conclude that $$\text{[Dirichlet plate w/hole]}\text{[Complementary Neumann disk]}=\text{[Entire Dirichlet plane]}$$ (6) and similarly $$\text{[Neumann plate w/hole]}\text{[Complementary Dirichlet disk]}=\text{[Entire Neumann plane]}.$$ (7) ## IV Line Segment in Two Spatial dimensions In this section we consider a scalar field in 2+1 dimensions with boundary conditions imposed on a line segment from $`x=d`$ to $`x=d`$. In circular coordinates, we can decompose a free, real, massless scalar field in modes as $$\varphi (r,\theta )=\underset{m=0}{\overset{\mathrm{}}{}}{}_{}{}^{}_{0}^{\mathrm{}}dk\sqrt{\frac{k}{2\pi \omega }}J_m(kr)(\mathrm{cos}m\theta b_k^m{}_{}{}^{}+\mathrm{sin}m\theta c_k^m{}_{}{}^{})e^{i\omega t}+\text{c.c.}$$ (8) where the prime on the summation sign indicates that for $`m=0`$, there is no sin mode and instead of $`\mathrm{cos}0=1`$ we have $`1/\sqrt{2}`$. Now we go to elliptical coordinates. For notational consistency with the three-dimensional case, we consider the $`x`$-$`z`$ plane. The foci will be located at $`x=\pm d`$, and the the geodesic will run along the $`z`$ axis. Elliptical coordinates $`\mu `$, $`\theta `$ are given by $`x`$ $`=`$ $`d\mathrm{cosh}\mu \mathrm{cos}\theta `$ (9) $`z`$ $`=`$ $`d\mathrm{sinh}\mu \mathrm{sin}\theta `$ (10) so $$r=\sqrt{x^2+z^2}=d\sqrt{\frac{\mathrm{cosh}2\mu +\mathrm{cos}2\theta }{2}}\frac{d}{2}e^\mu $$ (11) as $`\mu \mathrm{}`$. We define our Mathieu functions following the conventions of Abramowitz and Stegun Abramowitz and Stegun (1972) but extending their notation to be more similar to that of Bessel functions. The angular functions are $`ce_m(\theta ,q)`$ and $`se_m(\theta ,q)`$. They satisfy $$y^{\prime \prime }+(a2q\mathrm{cos}2\theta )y=0$$ (12) where $`q=(dk/2)^2`$, and are normalized so that $$_0^{2\pi }𝑑\theta ce_m(\theta ,q)^2=_0^{2\pi }𝑑\theta se_m(\theta ,q)^2=\pi .$$ (13) This normalization holds even for $`ce_0`$, but there is no such function as $`se_0`$. Thus the normalization is precisely the same as the circular functions used above, including the special case for $`m=0`$. The radial functions of the first kind are $`Je_m(\mu ,q)`$ and $`Jo_m(\mu ,q)`$ and are precisely the functions called $`Mc_m^{(1)}(\mu ,q)`$ and $`Ms_m^{(1)}(\mu ,q)`$ respectively in Abramowitz and Stegun (1972). They satisfy $$y^{\prime \prime }(a2q\mathrm{cosh}2\mu )y=0$$ (14) and go asymptotically to $`J_m(\sqrt{q}e^\mu )=J_m(kr)`$. Note that the functions $`Je_m`$ and $`Jo_m`$ defined in Morse and Feshbach (1953) have an additional factor of $`\sqrt{\pi /2}`$. Analogously, we denote the radial functions of the second kind as $`Ye_m(\mu ,q)`$ and $`Yo_m(\mu ,q)`$, which are $`Mc_m^{(2)}(\mu ,q)`$ and $`Ms_m^{(2)}(\mu ,q)`$ in Abramowitz and Stegun (1972). The normalization of radial functions depends only on their asymptotics, so $`Je_m`$ and $`Jo_m`$ have the same normalization as $`J_m`$. Thus the field becomes $$\varphi (\mu ,\theta )=\underset{m=0}{\overset{\mathrm{}}{}}{}_{}{}^{}_{0}^{\mathrm{}}dk\sqrt{\frac{k}{2\pi \omega }}(Je_m(\mu ,q)ce_m(\theta ,q)b_k^m{}_{}{}^{}+Jo_m(\mu ,q)se_m(\theta ,q)c_k^m{}_{}{}^{})e^{i\omega t}+\text{c.c.}$$ (15) where the prime on the summation sign indicates that the second term is included only for $`m>0`$. Now we consider Neumann conditions along the line segment, which is $`\mu =0`$. The even wavefunctions obey the conditions already, because $`dJe_m/d\mu =0`$ at $`\mu =0`$, but the odd functions need to be modified. Instead of $`\psi o_k^m=0`$ we need $`d\psi o_k^m/d\mu =0`$. For the free case we had $$Jo_m(\mu ,q)=\frac{1}{2}\left[Ho_m^{(1)}(\mu ,q)+Ho_m^{(2)}(\mu ,q)\right]$$ (16) where $`Ho^{(1)}=Jo_m+iYo_m`$ is the function called $`Ms_m^{(3)}`$ in Abramowitz and Stegun (1972) and $`Ho_m^{(2)}=Jo_miYo_m`$ is $`Ms_m^{(4)}`$. Now we need $$\psi o_m(\mu ,q)=\frac{1}{2}\left[e^{2i\delta }Ho_m^{(1)}(\mu ,q)+Ho_m^{(2)}(\mu ,q)\right]$$ (17) where $$e^{2i\delta }=\frac{Ho_m^{(2)}{}_{}{}^{}(0,q)}{Ho_m^{(1)}{}_{}{}^{}(0,q)}$$ (18) and the derivative is with respect to $`\mu `$. We can then compute the renormalized vacuum expectation value of the time-derivative term in Eq. (5), $`\dot{\varphi }^2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}𝑑kk\omega \left(\|\psi o_m(\mu ,q)\|^2Jo_m(\mu ,q)^2\right)se_m(\mu ,q)^2`$ (19) and we can write $$\|\psi o_m(\mu ,q)\|^2Jo_m(\mu ,q)^2=\frac{1}{4}\left[\left(e^{2i\delta }1\right)Ho_m^{(1)}(\mu ,q)^2+\left(e^{2i\delta }1\right)Ho_m^{(2)}(\mu ,q)^2\right].$$ (20) We want to extend the range of integration to include negative $`k`$. The term $`se_m(\mu ,q)^2`$ is unchanged by going to negative $`k`$, while $`Jo`$ and $`Ho`$ behave just like the corresponding Bessel functions. Thus the situation is exactly as in the circular case Saharian (2001); Schwartz-Perlov and Olum (2003): extending the range of integration exchanges the two terms in Eq. (20), so we can consider only the first term integrated over the entire real axis. When we close the contour at infinity, we get the contribution from the branch cut in $`\omega =\sqrt{k}`$. With $`k=i\kappa `$, the angular function becomes $`se_m(\theta ,q)`$, and the radial functions have the same continuation as Bessel functions, $`Jo_m(\mu ,q)`$ $`=`$ $`i^mIo_m(\mu ,q)`$ (21) $`Ho_m^{(1)}(\mu ,q)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}i^{(m+1)}Ko_m(\mu ,q)`$ (22) where $`Io_m`$ and $`Ko_m`$ are exactly as in Abramowitz and Stegun (1972). Thus $$e^{2i\delta _m(i\kappa )}1=()^{m+1}i\pi \frac{Io_m^{}(0,\phi )}{Ko_m^{}(0,\phi )}$$ (23) where $`\phi =(d\kappa /2)^2=q.`$ Putting it all together we have $`\dot{\varphi }^2`$ $`=`$ $`{\displaystyle \frac{()^mi}{\pi ^3}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}𝑑\kappa \left(e^{2i\delta }1\right)\kappa ^2Ko_m(\mu ,\phi )^2se_m(\theta ,\phi )^2`$ (24) $`=`$ $`{\displaystyle \frac{1}{\pi ^2}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}𝑑\kappa {\displaystyle \frac{Io_m^{}(0,\phi )}{Ko_m^{}(0,\phi )}}\kappa ^2Ko_m(\mu ,\phi )^2se_m(\theta ,\phi )^2.`$ (25) On the axis, terms with $`m`$ even vanish, so we have $$\dot{\varphi }^2=\frac{1}{\pi ^2}\underset{m=1}{\overset{\mathrm{}}{}}{}_{}{}^{}_{0}^{\mathrm{}}𝑑\kappa \frac{Io_m^{}(0,\phi )}{Ko_m^{}(0,\phi )}\kappa ^2Ko_m(\mu ,\phi )^2se_m(\pi /2,\phi )^2$$ (26) where the prime on the summation sign indicates that we sum over odd values of $`m`$. The other vacuum expectation value that we need is $`(_z\varphi )^2`$. On the $`z`$ axis, $`_z\varphi `$ is just the component of the gradient in the $`\mu `$ direction, which differs from $`_\mu \varphi `$ by the inverse of the metric coefficient $$h=d\sqrt{\mathrm{cosh}^2\mu \mathrm{cos}^2\theta }=d\sqrt{\frac{\mathrm{cosh}2\mu \mathrm{cos}2\theta }{2}}.$$ (27) The calculation is otherwise similar. Instead of two powers of $`\omega `$ from time differentiation we just have the radial function differentiated with respect to $`\mu `$, $$(_z\varphi )^2=\frac{1}{h^2}(_\mu \varphi )^2=\frac{1}{\pi ^2h^2}\underset{m=1}{\overset{\mathrm{}}{}}{}_{}{}^{}_{0}^{\mathrm{}}𝑑\kappa \frac{Io_m^{}(0,\phi )}{Ko_m^{}(0,\phi )}Ko_m^{}(\mu ,\phi )^2se_m(\pi /2,\phi )^2.$$ (28) If instead we have Dirichlet conditions on the line segment, the odd functions will be unmodified, but for the even functions we need $`\psi e_k^m=0`$ at $`\mu =0`$, so we have $$\psi e_m(\mu ,q)=\frac{1}{2}\left[e^{2i\delta }He_m^{(1)}(\mu ,q)+He_m^{(2)}(\mu ,q)\right]$$ (29) with $$e^{2i\delta }=\frac{He_m^{(2)}(0,q)}{He_m^{(1)}(0,q)}$$ (30) so $$e^{2i\delta _m(i\kappa )}1=()^{m+1}i\pi \frac{Ie_m(0,\phi )}{Ke_m(0,\phi )}.$$ (31) Thus on the axis we have $`\dot{\varphi }^2`$ $`=`$ $`{\displaystyle \frac{1}{\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle {}_{}{}^{}_{0}^{\mathrm{}}}𝑑\kappa {\displaystyle \frac{Ie_m(0,\phi )}{Ke_m(0,\phi )}}\kappa ^2Ke_m(\mu ,\phi )^2ce_m(\pi /2,\phi )^2`$ (32) $`(_z\varphi )^2`$ $`=`$ $`{\displaystyle \frac{1}{\pi ^2h^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle {}_{}{}^{}_{0}^{\mathrm{}}}𝑑\kappa {\displaystyle \frac{Ie_m(0,\phi )}{Ke_m(0,\phi )}}Ke_m^{}(\mu ,\phi )^2ce_m(\pi /2,\phi )^2.`$ (33) where the star on the summation sign indicates that we sum over even values of $`m`$. ## V Disk in Three Spatial Dimensions In this section we consider a scalar field with boundary conditions imposed on a disk of radius $`d`$ in the $`x`$-$`y`$ plane, centered at the origin. In spherical coordinates, we can decompose a free, real, massless scalar field in modes as $$\varphi (r,\theta ,\varphi )=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}_0^{\mathrm{}}𝑑k\frac{k}{\sqrt{\pi \omega }}j_l(kr)Y_{lm}(\theta ,\varphi )e^{i\omega t}+\text{c.c.}$$ (34) where $`j_l`$ is the spherical Bessel function. Next we go to oblate spheroidal coordinates, given by $`x`$ $`=`$ $`d\sqrt{(\xi ^2+1)(1\eta ^2)}\mathrm{cos}\varphi =d\mathrm{cosh}\mu \mathrm{sin}\theta \mathrm{cos}\varphi `$ (35) $`y`$ $`=`$ $`d\sqrt{(\xi ^2+1)(1\eta ^2)}\mathrm{sin}\varphi =d\mathrm{cosh}\mu \mathrm{sin}\theta \mathrm{sin}\varphi `$ (36) $`z`$ $`=`$ $`d\eta \xi =d\mathrm{sinh}\mu \mathrm{cos}\theta `$ (37) where $`\varphi `$ is the azimuthal angle, $`\eta =\mathrm{cos}\theta `$ is the coordinate akin to polar angle, and $`\xi =\mathrm{sinh}\mu `$ is the radial coordinate, with $`r=d\xi `$ for large $`\xi `$. We define the prolate angular spheroidal function $`S_n^m(c;\eta )=S_n^{m(1)}(c;\eta )`$ with $`c=kd`$, using the normalization of Meixner and Schäfke Meixner and Schäfke (1954). Prolate spheroidal functions can be converted to the oblate ones appropriate to our situation by $`kik`$ and $`\xi i\xi `$. Thus our angular functions are $`S_n^m(ic;\eta )`$, obeying the orthonormality relation $$_1^1S_n^m(ic;\eta )S_n^{}^m(ic;\eta )^2𝑑\eta =\frac{2}{2n+1}\frac{(n+m)!}{(nm)!}\delta _{nn^{}}$$ (38) Using these functions we can define oblate spheroidal harmonics by analogy with spherical harmonics, $$𝒴_n^m(ic;\eta ,\varphi )=\sqrt{\frac{2n+1}{4\pi }\frac{(nm)!}{(n+m)!}}S_n^m(ic;\eta )e^{im\varphi }$$ (39) obeying the analogous orthonormality relation $$_1^1𝑑\eta _0^{2\pi }𝑑\varphi 𝒴_n^m(ic;\eta ,\varphi )^{}𝒴_n^{}^m^{}(ic;\eta ,\varphi )=\delta _{nn^{}}\delta _{mm^{}}.$$ (40) We also define the radial spheroidal functions $`R_n^{m(1)}(ic;i\xi )`$ and $`R_n^{m(2)}(ic;i\xi )`$, normalized by $`\underset{\xi \mathrm{}}{lim}R_n^{m(1)}(ic;i\xi )`$ $`=`$ $`j_n(c\xi )`$ (41) $`\underset{\xi \mathrm{}}{lim}R_n^{m(2)}(ic;i\xi )`$ $`=`$ $`y_n(c\xi ).`$ (42) The radial functions thus have the same normalization as spherical Bessel functions. Thus the field becomes $$\varphi (\xi ,\eta ,\varphi )=\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=n}{\overset{n}{}}_0^{\mathrm{}}𝑑k\frac{k}{\sqrt{\pi \omega }}𝒴_n^m(ic;\eta ,\varphi )R_n^{m(1)}(ic;i\xi )e^{i\omega t}+\text{c.c.}$$ (43) If $`m+n`$ is even, then $`S_n^m(ic;\eta )`$ is an even function of $`\eta `$ and $`(d/d\xi )R_n^{m(1)}(ic;i\xi )=0`$ at $`\xi =0`$. Thus such wave functions will be continuous across the disk $`\eta =0`$. Similarly, if $`m+n`$ is odd, then $`S_n^m(ic;\eta )`$ is an odd function of $`\eta `$ and $`R_n^{m(1)}(ic;0)=0`$, so the product is once again continuous. The $`R^{(2)}`$ functions do not have these boundary conditions, so they cannot be used in the vacuum wave functions. Now we consider Neumann conditions on the disk $`\xi =0`$. If $`m+n`$ is even, the functions obey the conditions already. Otherwise we need to combine $`R_n^{m(1)}`$ and $`R_n^{m(2)}`$ to give the desired condition. With $`R_n^{m(3)}=R_n^{m(1)}+iR_n^{m(2)}`$ and $`R_n^{m(4)}=R_n^{m(1)}iR_n^{m(4)}`$ we can write the desired radial function $$\psi _n^m(ic;i\xi )=\frac{1}{2}\left[e^{2i\delta (ic)}R_n^{m(3)}(ic;i\xi )+R_n^{m(4)}(ic;i\xi )\right]$$ (44) with the condition $$e^{2i\delta (ic)}=\frac{R_n^{m(4)}{}_{}{}^{}(ic;0)}{R_n^{m(3)}{}_{}{}^{}(ic;0)}$$ (45) where the derivative is with respect to the second argument. The vacuum expectation value of the time derivative term, subtracting the free vacuum, then becomes $`\dot{\varphi }^2`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=n}{\overset{n}{}}}{\displaystyle {}_{}{}^{}_{0}^{\mathrm{}}}𝑑kk^2\omega \|𝒴_n^m(ic;\eta ,\varphi )\|^2\left[\|\psi _m(ic;i\xi )\|^2R_n^{m(1)}(ic;i\xi )^2\right]`$ (46) $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=n}{\overset{n}{}}}{\displaystyle {}_{}{}^{}_{0}^{\mathrm{}}}dkk^2\omega \|𝒴_n^m(ic;\eta ,\varphi )\|^2[(e^{2i\delta (ic)}1)R_n^{m(3)}(ic;i\xi )^2`$ (48) $`+(e^{2i\delta (ic)}1)R_n^{m(4)}(ic;i\xi )^2]`$ where the prime on the summation sign means that only odd values of $`m+n`$ are included. We want to extend the range of integration to include negative $`k`$, which changes the sign of $`c`$. We can implement this change by changing the sign of $`\xi `$ in $`R_n^{m(1)}`$ and $`R_n^{m(2)}`$, since these real functions are not affected by complex conjugation. The functions $`R_n^{m(1)}`$ and $`R_n^{m(2)}`$ have opposite parity under this transformation, so $`R_n^{m(3)}`$ and $`R_n^{m(4)}`$ change places. Thus including negative $`k`$ in the first term gives the second term, and we have $$\dot{\varphi }^2=\frac{1}{4\pi }\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=n}{\overset{n}{}}{}_{}{}^{}_{\mathrm{}}^{\mathrm{}}𝑑kk^2\omega \|𝒴_n^m(ic;\eta ,\varphi )\|^2\left(e^{2i\delta (ic)}1\right)R_n^{m(3)}(ic;i\xi )^2.$$ (49) If we take $`k`$ and thus $`c`$ in the upper half plane we will get spheroidal functions whose parameter goes to negative real infinity. Therefore we can close the contour at infinity and obtain an integral along the branch cut on the imaginary axis associated with the square root in $`\omega `$, $`\dot{\varphi }^2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=n}{\overset{n}{}}}{\displaystyle {}_{}{}^{}_{0}^{\mathrm{}}}𝑑\kappa \kappa ^3\|𝒴_n^m(\gamma ;\eta ,\varphi )\|^2\left(e^{2i\delta (\gamma )}1\right)R_n^{m(3)}(\gamma ;i\xi )^2`$ (50) $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=n}{\overset{n}{}}}{\displaystyle {}_{}{}^{}_{0}^{\mathrm{}}}𝑑\kappa \kappa ^3{\displaystyle \frac{R_n^{m(1)}{}_{}{}^{}(\gamma ;0)}{R_n^{m(3)}{}_{}{}^{}(\gamma ;0)}}\|𝒴_n^m(\gamma ;\eta ,\varphi )\|^2R_n^{m(3)}(\gamma ;i\xi )^2`$ (51) where $`\gamma =ic=ikd=\kappa d`$. On the axis, we have $$\dot{\varphi }^2=\frac{1}{\pi }\underset{n=1}{\overset{\mathrm{}}{}}{}_{}{}^{}_{0}^{\mathrm{}}𝑑\kappa \kappa ^3\frac{R_n^{0(1)}{}_{}{}^{}(\gamma ;0)}{R_n^{0(3)}{}_{}{}^{}(\gamma ;0)}\|𝒴_n^0(\gamma ;1,\varphi )\|^2R_n^{0(3)}(\gamma ;i\xi )^2$$ (52) where we have specialized to $`m=0`$ because the contributions from nonzero $`m`$ vanish on the axis, leaving only a sum over odd values of $`n`$. Similarly, on the axis we have $$(_z\varphi )^2=\frac{1}{\pi d^2}\underset{n=1}{\overset{\mathrm{}}{}}{}_{}{}^{}_{0}^{\mathrm{}}d\kappa \kappa \frac{R_n^{0(1)}{}_{}{}^{}(\gamma ;0)}{R_n^{0(3)}{}_{}{}^{}(\gamma ;0)}\|𝒴_n^0(\gamma ;1,\varphi )\|^2R_n^{0(3)}{}_{}{}^{}(\gamma ;i\xi )_{}^{2}$$ (53) where the primes on the radial functions indicate derivatives with respect to the second argument, and the metric coefficient $$h_\xi =\left\|\frac{𝐱}{\xi }\right\|=d\sqrt{\frac{\xi ^2+\eta ^2}{\xi ^2+1}}.$$ (54) becomes equal to $`d`$ on the axis. For Dirichlet conditions on the disk, we have the analogous results $`\dot{\varphi }^2`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle {}_{}{}^{}_{0}^{\mathrm{}}}𝑑\kappa \kappa ^3{\displaystyle \frac{R_n^{0(1)}(\gamma ;0)}{R_n^{0(3)}(\gamma ;0)}}\|𝒴_n^0(\gamma ;\eta ,\varphi )\|^2R_n^{0(3)}(\gamma ;i\xi )^2`$ (55) $`(_z\varphi )^2`$ $`=`$ $`{\displaystyle \frac{1}{\pi d^2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle {}_{}{}^{}_{0}^{\mathrm{}}}d\kappa \kappa {\displaystyle \frac{R_n^{0(1)}(\gamma ;0)}{R_n^{0(3)}(\gamma ;0)}}\|𝒴_n^0(\gamma ;1,\varphi )\|^2R_n^{0(3)}{}_{}{}^{}(\gamma ;i\xi )_{}^{2}`$ (56) where we have again specialized to the axis so that only $`m=0`$ contributes, the derivatives of the radial functions are again with respect to the second argument, and the star on the summation sign indicates that now we sum over even values of $`n`$. ## VI Numerical Calculation For a null geodesic $`V^\lambda `$ perpendicular to the plate and passing through the center of the hole, the contribution to ANEC is given by Eq. (5). We can compute the results for the complementary disk using the formulae derived in the previous sections. We then add the complete plate results, $$\dot{\varphi }^2+(_z\varphi )^2=\{\begin{array}{cc}\frac{1}{32\pi z^3}\hfill & \text{in two dimensions, and}\hfill \\ \frac{1}{16\pi ^2z^4}\hfill & \text{in three dimensions}\hfill \end{array}$$ (57) with Dirichlet conditions giving the upper sign and Neumann the lower. In two dimensions, we compute the Mathieu functions using the package of Alhargan Alhargan (2000a, b), with some minor modifications: we use 80-bit double precision throughout the calculation to accommodate the extreme dynamical range needed for the wide range of Mathieu function parameters we use, and we have adapted the code to use our set of normalization conventions. These C++ routines are then imported into Mathematica, where we can use efficient routines for numerical sums and integrals. In three dimensions, we use the Mathematica spheroidal harmonic package of Falloon Falloon et al. (2003). We have updated it to fix incompatibilities with the latest version of Mathematica and to avoid memory leaks. We have also made a number of efficiency optimizations appropriate to the unusual demands we make on the code (for example, we changed the caching structure so that it is appropriate to the way we call the functions, with the same arguments but different parameters rather than the other way around; we also wrote specific code for the modified radial function of the third kind to avoid cancellations of exponentially growing quantities). Figure 3 shows the contributions to NEC for Dirichlet plates with holes of unit radius in two and three spatial dimensions, as functions of distance along the axis. Using Babinet’s principle, we have computed the sum of contributions from a Neumann disk and an infinite Dirichlet mirror. At small distances, the contributions from both the finite disk and the infinite mirror diverge like $`1/z^{n+1}`$, where $`z`$ is the distance from the origin and $`n`$ is the spatial dimension. The true result, however, does not diverge (the origin is just a point in empty space) and by symmetry must have zero slope at the origin. This cancellation provides a highly nontrivial check on our calculation. Going all the way to the origin would require infinite precision; in two dimensions we stop at a distance $`0.15d`$, while in three dimensions we stop at distance $`0.25d`$. At these values, our curves already show this cancellation clearly. We also extrapolate our result (without putting in any restrictions on the extrapolation at the origin) and find that it goes smoothly to a finite value with zero slope at $`z=0`$. A less stringent check is that our calculation approaches the perfect mirror at large distances; in our approach this result simply tells us that the finite disk contribution is going to zero fast enough. ## VII Discussion The results in Fig. 3 are quite striking. In both cases, far from the origin we see the negative contribution to ANEC that we would expect from the standard calculation. Near the plate, however, the hole leads to a large positive contribution. Integrating the results shown in Fig. 3 gives a total contribution (including both sides of the plate) of $`1.63\times 10^3/d^3`$ in three dimensions and $`4.53\times 10^3/d^2`$ in two dimensions. The positive contribution overwhelms the negative contribution, so that ANEC is obeyed. We can look at these results from a different point of view by considering conformal coupling. Since the NEC contribution in this case differs from that in minimal coupling by a total derivative, it leads to the same results for ANEC. From this point of view, one might not expect any ANEC violation in the case at hand, because the quantum contributions to the perfect mirror vanish. (In the perfect mirror case, changing from minimal to conformal coupling effectively moves the negative contributions from the region outside the boundary onto the boundary itself Ford and Svaiter (1998).) However, from this point of view one might just as well expect Neumann conditions to violate ANEC; while the perfect Neumann mirror result is positive in minimal coupling, it also vanishes for conformal coupling. The results for Neumann conditions are shown in Fig. 4. Once again integration gives positive results, $`2.10\times 10^3/d^3`$ in three dimensions and $`1.77\times 10^2/d^2`$ in two dimensions, so ANEC is obeyed. Thus again from this point of view, one finds ANEC obeyed more often than would be naïvely expected. A conformal field between two plates would have a constant negative energy density. We can estimate how our results would extend to the case of parallel plates with holes, in two limits. If the separation between the plates $`\mathrm{}`$ is much smaller than the radius of the hole $`d`$, the two plates are equivalent to a single plate and ANEC continues to hold. In the other extreme, if the separation between the plates is large compared to the radius of the hole, then we can assume that the change $`\mathrm{\Delta }`$ in ANEC induced by adding a hole in one plate is unaffected by the other plate. In three dimensions, we obtain a contribution to ANEC of $$𝑑xV^\lambda V^\nu T_{\lambda \nu }2\mathrm{\Delta }\frac{\pi ^2}{720\mathrm{}^3}$$ (58) where the first term is the effect of the each of the two holes individually and the second term is the standard contribution from the two plates. For Dirichlet conditions, $`\mathrm{\Delta }=1.63\times 10^3/d^3`$. Thus Eq. (58) gives a positive result as long as $`\mathrm{}>1.6d`$, which surely includes its entire range of applicability. Similar results hold for the other cases we have considered. ## VIII Acknowledgments K. D. O. thanks Xavier Siemens for helpful conversations. N. G. thanks the Kavli Institute for Theoretical Physics for hospitality while part of this work was completed. K. D. O. was supported in part by the National Science Foundation (NSF). N. G. was supported in part by the NSF through the Vermont Experimental Program to Stimulate Competitive Research (VT-EPSCoR).
warning/0506/hep-ph0506161.html	ar5iv	text	# Saturation 2005 ( mini-review) ## I Saturation at HERA and RHIC Before discussing the high density QCD news we would like to summarize what we have learned about saturation at HERA and RHIC. HERA: * The power - like growth of $`xG(x,Q^2)`$ at low $`x`$ ($`xG(x,Q^2)x^\lambda `$ with $`\lambda 0.3`$; * The geometrical scaling behaviour for $`x10^2`$; * Fit of all HERA data for $`Q^2=0÷500GeV^2`$ with $`\chi ^2/d.o.f.\mathrm{\hspace{0.17em}1}`$ based on non-linear equation Levin:GLLM ; Levin:IIM ; RHIC: * Saturation approach for $`dN/dy`$ versus $`y`$, energy and number of participants predicted and led to a reasonable description of the experimental data Levin:KLN ; * Prediction for suppression of the hadron production in dA collision and confirmation in the experimental data LEVIN:KLM ; Levin:KKT . The only consistent explanation all these observations is to assume that at HERA we have started to approach a new phase of QCD, with large gluon density but still with small coupling constant. The regime of high parton density at HERA is reached due to the QCD emission of gluons that was incorporated in the QCD evolution equations. The independent check of the effects of high gluon density at HERA was performed by RHIC experiment in heavy ion-ion collisions. In this reaction the energies are much lower than at HERA, but the large values of the parton densities were achieved due to the large number of nucleons in a nucleus. Based on these experimental observations we can anticipate that the LHC will be a machine for discovery a new phase of QCD: colour glass condensate with saturated gluon density. ## II Predictions for the LHC range of energies Our main challenge is to provide reliable estimates for the influence of high density QCD (saturation) effects in the LHC range of energies. The first such estimates have been discussed Levin:GLMN ; Levin:ESQU , and the results for the ratio of the unintegrated structure functions $`D=\varphi ^{NL}/\varphi ^L`$ are plotted in Fig.1 where $$\frac{d\sigma }{dyd^2p_t}\frac{\alpha _S}{p_t^2}d^2k_t\varphi (k_t^2)\varphi ((\stackrel{}{p}\stackrel{}{k})_t^2)$$ (1) and $`\varphi ^{NL}(\varphi ^L)`$ is solution of the non-linear (linear ) equation. It should be stressed that non-linear evolution predicts not only suppression in the saturation region, but also the anti-shadowing effect which results in an increase of the value of $`\varphi `$ for $`Q^2>Q_s^2(x)`$, where $`Q_s`$ is the saturation scale. One can see that the suppression and increase could be rather large leading to an inclusive cross section twice as large or twice as small, as the predictions based on routine linear evolution. ## III Theoretical development ### III.1 B- JIMWLK approach $``$ BFKL Pomeron Calculus The good news is that it turns out that Balitsky-JIMWLK approach Levin:JIMWLK can be reduced to BFKL Pomeron calculus Levin:KL , and JIMWLK effective Lagrangian give us possibility to calculate all multi-Pomeron vertices. For the first time, we can do such calculations using operator formalism without spending years to obtain result just summing Feyman diagrams. Since the colour dipoles are the ‘wee’ partons of the BFKL equation the Balitsky-JIMWLK formalism can be discussed in terms of the dipole approach. The bad news is that we have not achieved any progress in Pomeron calculus. ### III.2 Probabilistic interpretation Our last hope is the probabilistic approach to Pomeron interaction. The best way to express our optimism is to cite Grassberger and Sundermeyer Levin:GS who proposed this interpretation: “ Reggeon field theory is equivalent to a chemical process where a radical can undergo diffusion, absorption, recombination, and autocatalytic production. Physically, these ”radicals” are wee partons (colour dipoles)”. It turns out that B-WLKJIM approach can be written as a typical death-birth process (Markov’s chain)Levin:BIW ; Levin:LL $$\frac{P_n}{Y}=\underset{i}{}\mathrm{\Gamma }(12)\left(P_n(\mathrm{}x_i,y_i\mathrm{};Y)P_{n1}(\mathrm{}x_i,y_i\mathrm{};Y)\right)$$ (2) where $`P_n`$ \- probability to find $`n`$-dipoles at rapidity $`Y`$, $`\mathrm{\Gamma }(12)`$ describe the decay of one dipole into two dipoles and $``$ denotes all needed integration. This equation can be a basis for the Monte Carlo code which will be able to solve high density QCD equations, and which will lead to theoretical treatment of the multiparticle production. ### III.3 Hunt for Pomeron loops The process of two Pomeron merging into one Pomeron is naturally included in Pomeron calculus with the same vertex as the process of Pomeron splitting. However, we need correctly normalize this process if we wish to use the probabilistic interpretation. Such normalization was suggested in Ref. Levin:IT and this vertex $`\mathrm{\Gamma }(21)`$ has been calculated Levin:IT ; Levin:MSW ; Levin:LL . Using this vertex, we can generalize Eq.(2) which takes the form $$\frac{P_n}{Y}=\text{Eq.(2)}\underset{i}{}\mathrm{\Gamma }(21)\left(P_n(\mathrm{}x_i,y_i\mathrm{};Y)\underset{k}{}P_{n+1}(\mathrm{}x_i,y_i\mathrm{}x_k,y_k;Y)\right)$$ (3) ### III.4 Solution Attempts to solve Eqs.(3) have been made in Refs.Levin:BO ; Levin:L ; Levin:RS . The result is surprisingly unexpected, namely, * Asymptotic solution leads to a gray disc (not black!!!); * Using the large parameters of our theory ($`\mathrm{\Gamma }(12)/\mathrm{\Gamma }(21)N_c^2/\alpha _S^2`$ and $`\mathrm{\Gamma }(12)/\mathrm{\Gamma }(23)N_c^2`$) the semiclassical approach can be developed for searching for both the asymptotic solution and the corrections to it, at high energy; * The corrections to the asymptotic solution decrease at large values of $`Y`$, and can be found from the Liouville-type linear equation; * The important role in searching for high energy asymptotic behaviour of the amplitude plays the role of $`t`$-channel unitarity constraint, which specifies the value of the typical amplitude for dipole-dipole interaction. ### III.5 Topics which I have no room to discuss This brief review is my personal view on news in low $`x`$ (high density) QCD. Unfortunately, I had no room even to express my point of view. It is pity since I think that a more microscopic approach, related to the new effective Lagrangian, and to a search for a Bogolubov transformation between dipole and quarks (antiquark) and gluon degrees of freedom Levin:KL ; Levin:HIMST ; Levin:MMSW , looks very interesting. It is very attractive approach and I hope that my references provide the reader with names of active players in this field. However, I must admit that the theory becomes dangerously complicated and reminds me more and more my nightmare that Lipatov Levin:LI is correct with his effective action, which is not easier to solve than the full QCD Lagrangian. Acknowledgments: I am very grateful to E. Gotsman for everyday discussions on the subject of this talk.
warning/0506/astro-ph0506570.html	ar5iv	text	# Non-Gaussianity from Cosmic Magnetic Fields ## I Introduction It appears increasingly likely from cosmological data such as the cosmic microwave background (CMB) that the dominant source of perturbations in the universe was adiabatic and Gaussian in nature, consistent with the expectations of an inflationary universe Bennett et al. (2003); Komatsu et al. (2003). However, this does not exclude the possibility that non-linear sources might also have played a role in sourcing perturbations. While their impact on the power spectrum might be small, they could dominate any non-Gaussianities that we observe. In order to search for these optimally we need clear predictions for the nature of the non-Gaussian signals that they would source. Here we begin to investigate the signatures for a particular non-linear model, early-universe magnetic fields. Magnetic fields are observed on many scales in the cosmos, from planetary scales with fields of strengths of a few Gauss, to galactic scales at a strength of approximately $`\mu G`$. They also exist between galaxies and are likely to be in clusters, both with field strengths somewhere between the nano- and micro-Gauss level and a coherence length on the order of megaparsecs. While fields on supercluster scales are extraordinarily difficult to detect, there are suggestions that fields up to the order of $`\mu G`$ may exist even there. (See for example Kronberg (1994); Kim et al. (1991); Zweibel and Heiles (1997); Grasso and Rubenstein (2001); Widrow (2002); Giovannini (2004a) for reviews.) The origin of these fields is uncertain. Possible creation mechanisms can be separated into processes occurring before, during or after recombination. Suggested post-recombination processes include the dynamo mechanism Grasso and Rubenstein (2001); Zeldovich et al. (1980) (wherein magnetic energy is provided by rotational kinetic energy), or the adiabatic compression of an already-magnetized cloud Grasso and Rubenstein (2001). However, these still rely on a seed field, possibly created by some large battery mechanism; an efficient dynamo could amplify a field of strength $`10^{30}G`$ to observed levels, while an adiabatic compression requires a seed many orders of magnitude larger. While mechanisms certainly exist to generate this field at reionization Gnedin et al. (2000); Langer et al. (2003) or at recombination itself Hogan (2000); Berezhiani and Dolgov (2004), they might also have been created before recombination or even before nucleosynthesis; such early seed fields are constrained by limits from cosmic microwave background (see Barrow et al. (1997); Clarkson et al. (2003); Banerjee and Jedamzik (2004)) and nucleosynthesis Grasso and Rubenstein (2001); Caprini and Durrer (2002). Here we concentrate on fields from the early universe, and there is no shortage of suggestions for how such fields might have arisen. They could have been produced directly by inflation, or might have been generated during the electroweak symmetry breaking phase Turner and Widrow (1988); Ratra (1992); Grasso and Rubenstein (2001); Dimopoulos et al. (2002); Giovannini (2004a); Prokopec and Puchwein (2004); Bamba and Yokoyama (2004); Bassett et al. (2001); Ashoorioon and Mann (2005). There are also recent studies into generic phase transitions generating large-scale fields (e.g., Boyanovsky et al. (2003); Boyanovsky and de Vega (2005)). Cosmic defects might also be responsible Dimopoulos (1998); Davis and Dimopoulos (2005). More recently, attention has been given to the possibility that the fields might have been created continuously in the period between lepton decoupling and recombination, through the vorticity naturally occurring at higher order in perturbation theory Betschart et al. (2004); Matarrese et al. (2005); Gopal and Sethi (2004); Takahashi et al. (2005). Even after they are produced, they could evolve in the very early universe, perhaps as a result of hydromagnetic turbulence (e.g., Brandenburg et al. (1996)). (Dolgov Dolgov (2003) provides a brief overview of many of these creation mechanisms.) Primordial magnetic fields can have a significant impact on the cosmic microwave background. While early treatments focussed on this effect for a Bianchi universe Jacobs (1969); Milaneschi and Fabbri (1985), more modern treatments Scannapieco and Ferreira (1997); Subramanian and Barrow (1998); Durrer et al. (1998); Koh and Lee (2000); Durrer et al. (2000); Subramanian and Barrow (2002); Seshadri and Subramanian (2001); Mack et al. (2002); Caprini et al. (2004); Subramanian et al. (2003); Berera et al. (2004); Lewis (2004); Giovannini (2004b); Yamazaki et al. (2004) consider either small perturbations around a large-scale homogeneous field or a ‘tangled’ field configuration. Both of these scenarios source CMB perturbations, either directly through the scalar, vector and tensor stresses, or indirectly by the density and velocity perturbations they induce in the charged proton-electron fluids. Code for calculating the vector and tensor anisotropies generated by primordial magnetic fields was recently added to the publicly-available CAMB Lewis et al. (2000); Lewis (2004) and that for scalars has been modelled independently Koh and Lee (2000); Giovannini (2004b); Yamazaki et al. (2004). In addition to creating CMB anisotropies, magnetic fields can affect also CMB polarization by inducing Faraday rotation Loeb and Kosowsky (1996); Scoccola et al. (2004); Campanelli et al. (2004); Kosowsky et al. (2005). It would be useful to have a firm grasp on the statistical impact magnetic fields might have upon the microwave background. Our ultimate goal is to translate the non-Gaussianity in the magnetic stresses into a unique non-Gaussian signature in the microwave sky Brown and Crittenden . While there are many potential sources of CMB non-Gaussianity, including the non-linear evolution of the perturbations and gravitational lensing (for example, Bartolo et al. (2004); Lesgourgues et al. (2005); Bartolo et al. (2005)), there are also many ways a map may be non-Gaussian. The hope is that the different physical mechanisms will each have a characteristic non-Gaussian signature which can be searched for in the CMB sky. It is known that magnetic fields can introduce non-Gaussianity in many ways, including via Alfvén turbulence or secondary perturbations, and each of these can be searched for in the CMB Naselsky et al. (2004); Chen et al. (2004); Gopal and Sethi (2005). In addition, helical magnetic fields can produce parity-odd correlations in the CMB polarization field Kahniashvili and Ratra (2005). Here we focus on the non-Gaussianity of the magnetic sources themselves, with the three point moment as the starting point for our investigations. To this end we model the statistics of Gaussian-random magnetic fields, concentrating on the two- and three-point statistics of the magnetic stresses and the cross-correlations between the scalar, vector and tensor components, usually considered independently. We do so by both generating realizations of Gaussian-random magnetic fields and through numerically integrating pure analytical results, most of which are presented here for the first time. Significant non-Gaussianities and cross-correlations are expected due to the quadratic nature of the magnetic stress tensor; this suspicion is confirmed and the bispectra are presented. In section II we review the basic model of tangled primordial magnetic fields, their impact on the CMB and their statistics, detail our realizations of these fields and construct the stress-energy tensor. In section III we briefly consider the one-point statistics of the scalar pressures of the magnetic fields, evaluating the probability distribution function, the skewness and the kurtosis for both the isotropic and the anisotropic pressures. In section IV we review the two-point results (previously demonstrated in part Mack et al. (2002)) and present the power spectra of the separate components of the stress-energy tensor, displaying the excellent agreement between the theory and the simulations. In section V we move onto the three-point moments, deriving and presenting analytical results for their bispectra and good agreement with the simulated fields. We discuss the implications of our results in section VI. ## II Tangled magnetic fields Rather than considering a large-scale homogeneous magnetic field with a small inhomogeneous perturbation, we instead consider entirely inhomogeneous fields tangled on some length scale (see for example Mack et. al. Mack et al. (2002)). For simplicity (and to compare with the previous literature), we assume that the magnetic fields are random variables obeying a Gaussian probability distribution function; however, our realizations may be formulated more generally than this. Motivated by linear perturbation theory, we also take the fields to be effectively frozen, with their overall energy density decreasing as radiation ($`𝐁a^2`$, where $`a`$ is the scale factor.) Due to the extremely high conductivity of the early universe Grasso and Rubenstein (2001); Turner and Widrow (1988), we also assume that the electric components of the electromagnetic field vanish. Magnetic fields, being a non-linear source with a full anisotropic stress, will naturally source scalar, vector and tensor perturbations. For scalar perturbations one does not expect a significant effect on large scales. However, the contribution to the temperature auto-correlation on the microwave sky can begin to dominate at an $`l`$ of about $`1,000`$ (see for example Subramanian and Barrow (2002); Lewis (2004); Yamazaki et al. (2004)). This effect comes both from the impact of the magnetic energy and anisotropic pressure directly onto the spacetime geometry and from the Lorentz forces imparted onto the coupled proton-electron fluid. One also expects a significant impact on the vector perturbations as compared to the standard picture since the magnetic fields will both directly generate vector perturbations in the spacetime and will also contribute a Lorentz force to the solenoidal component of the velocity. Prior to neutrino decoupling a primordial magnetic field acts as a source for gravitational waves; following neutrino decoupling there will also be a contribution from the neutrinos, which Lewis has shown serves to cancel much of the magnetic stress Lewis (2004). Since the Lorentz forces and the stress-energy tensor are both quadratic in the magnetic field, we also expect a level of non-Gaussianity to be imprinted onto the fluid and perturbations, even for a magnetic field that is itself Gaussian. The features of the magnetic field – the Lorentz forces and the direct sourcing of geometric fluctuations – are all contained within the stress-energy tensor. We can then consider the statistics of the stress-energy tensor alone and be hopeful of characterizing the majority of the non-Gaussian effects that might impact on the CMB. The non-Gaussianity predicted from our analysis and our simulations can be projected onto the CMB sky by folding them with the transfer functions generated by the modified CAMB code Lewis et al. (2000); Lewis (2004); while this does not include support for scalar perturbations it will give good predictions on all scales for the $`B`$-mode polarization correlations which is sourced purely by vector and tensor perturbations. However, care would have to be taken to detangle these from any $`B`$-modes caused by the gravitational lensing of the dominant $`E`$-mode polarization. <sup>1</sup><sup>1</sup>1As earlier noted, magnetic fields would also cause Faraday rotation from $`E`$ modes into $`B`$ modes Loeb and Kosowsky (1996); Scoccola et al. (2004); Campanelli et al. (2004); Kosowsky et al. (2005); support for scalar modes in a uniform magnetic field was added to CMBFast by Scoccola et al. (2004) and also to CMBFast for tangled fields by Kosowsky et al. (2005), both in the case of small rotations. ### II.1 Statistics of the fields We begin by specifying the statistics of the tangled magnetic fields. In Fourier space, the divergence free condition for magnetic fields implies that $$𝐁_a(𝐤)𝐁_b^{}(𝐤^{})=𝒫(k)P_{ab}(𝐤)\delta (𝐤𝐤^{})+\frac{i}{2}(k)ϵ_{abc}\widehat{k}_c$$ (2.1) where $`𝒫(k)`$ is the magnetic field power spectrum, $`P_{ab}`$ is the operator projecting vectors and tensors onto a plane orthogonal to $`\widehat{k}_a`$ and $`\widehat{k}_b`$, $$P_{ab}(𝐤)=\delta _{ab}\widehat{k}_a\widehat{k}_b,$$ (2.2) with $`\delta _{ab}=\mathrm{diag}(1,1,1)`$ and $`(k)`$ is the power spectrum of the anti-symmetric helical term (see for example Pogosian et al. (2002); Caprini et al. (2004); Kahniashvili and Ratra (2005).) Here we have assumed the fields are statistically isotropic and homogeneous. If the magnetic fields are Gaussianly distributed, then all their statistics are determined by their power spectrum. In the interests of simplicity we henceforth assume that the helical component of the field vanishes. In real space, the magnetic correlation function is the transform of the magnetic power spectrum, implying that $$𝐁_a(0)𝐁_b(𝐱)=\delta _{ab}C_0(x)+\frac{}{x_a}\frac{}{x_b}C_1(x)$$ (2.3) where $`C_0(x)=\frac{V}{(2\pi )^3}d^3𝐤𝒫(k)e^{i𝐤𝐱}`$ and $`C_1(x)=\frac{V}{(2\pi )^3}d^3𝐤(𝒫(k)/k^2)e^{i𝐤𝐱}`$. In the limit of very small separations, the correlation becomes $$𝐁_a(0)𝐁_b(𝐱)=\frac{2}{3}\delta _{ab}C_0(0)+x_ax_bC_2(0)$$ (2.4) where $`C_2(0)𝑑kk^4𝒫(k)`$ is another correlation function related to the power spectrum. Often in the literature, the power spectrum is taken to be a simple power law, $$𝒫(k)=Ak^n.$$ (2.5) To avoid divergences, these power law spectra are generally assumed to have some small scale cutoff associated with the photon viscosity damping scale; see Mack et. al. Mack et al. (2002) for an evaluation of these for vector and tensor perturbations generated by a stochastic field with a power-law spectrum. Durrer and Caprini Caprini and Durrer (2002); Durrer and Caprini (2003) demonstrate that, for a causally-generated magnetic field, the spectral index must be $`n2`$; it must also in any event be $`n3`$ to avoid over-production of long-range coherent fields. They also show that, for the tensor case at least, the anisotropic stress will have a white-noise ($`n_{\mathrm{eff}}=0`$) spectrum for all $`n3/2`$. They also demonstrate that nucleosynthesis limits on the gravitational waves produced by the magnetic fields place extremely strong bounds on magnetic fields, to the level of $`10^{39}G`$ for inflation-produced fields with $`n=0`$, although this has been contested (Kosowsky et al. (2005); Caprini and Durrer (2005)). For spectra that might realistically imprint on the microwave background we must consider spectra with $`n<2`$, which are much less tightly constrained. For purposes of comparison we shall usually take either $`n=0`$ for a flat spectrum – where “flat” refers to the spectrum of primordial magnetic fields themselves – in which each mode contributes equally, or $`n=2.5`$ for a spectrum nearing the “realistically observable” $`n3`$, which we shall refer to as “steep”. Such fields can be produced by inflation (for example, Ratra (1992); Bamba and Yokoyama (2004)). It is worth emphasising that, while we are restricting ourselves to a power-law spectrum, this is not a necessity. We are generating simulated fields and calculating the statistics from these; it is as simple to employ other power spectra as it is to employ a power law. It is also a simple matter to employ a non-sharp damping scale – we could, for example, employ an exponentially or Gaussian-damped tail above a particular scale $`k_D`$ and gain some greater freedom in modelling the effective microphysics. As a particular example, Matarrese et. al. Matarrese et al. (2005) derived the power spectrum they expect from a field sourced from second-order vorticity in the electron-baryon plasma. This power spectrum has no simple functional form but is instead presented as a numerical integration and, in principle, it would be simple to employ this integration as the input power spectrum. The normalization of the power spectrum is typically fixed by reference to a particular co-moving smoothing scale $`\lambda `$ and the variance of the field strength at this scale, $`B_\lambda `$. Specifically, we smooth the field by convolving it with the Gaussian filter $$f(k)=\mathrm{exp}\left(\frac{\lambda ^2k^2}{2}\right)$$ (2.6) and define the variance of the field strength at the scale $`\lambda `$ by $$B_\lambda ^2=B_a(𝐱)B_a(𝐱),$$ (2.7) implying that the power spectrum and $`B_\lambda `$ are related by $$B_\lambda ^2=d^3𝐤𝒫(k)e^{\lambda ^2k^2}.$$ (2.8) This allows us to relate the astronomically observed field strengths at, say, cluster scales, to the amplitude of the magnetic power spectrum. For simplicity and concreteness, we will assume throughout that the fields themselves are Gaussian (consistent with previous literature). Obviously, any conclusions about the non-Gaussian signatures of magnetic fields will depend sensitively on this assumption, and we plan to explore more general scenarios in future work. Another ansatz for the fields is that they possess a $`\chi ^2`$ probability distribution function; such fields might arise from weakly non-linear effects since they are sampled from the combination of two (first-order) Gaussian variables. For example, the sourcing of second-order vorticity in the electron-baryon plasma by first-order density perturbations could lead to a weak but continually generated magnetic field Betschart et al. (2004); Matarrese et al. (2005); Gopal and Sethi (2004); Takahashi et al. (2005). Magnetic fields sourced at recombination and later are unlikely to be reasonably described by the above formalism; those sourced at recombination Berezhiani and Dolgov (2004) for example naturally arise at a very small scale. The question of the statistical nature of both recombination and very early universe fields has not been well explored in the literature. Here we present results from a Gaussian field as an illustrative example, but further work is needed to explore the consequences of any given microphysical mechanism. ### II.2 Realizations of magnetic fields To aid our study of the non-Gaussian properties of tangled magnetic fields, we create static realizations of the fields numerically. We create the fields on a grid in Fourier space of size $`l_{\mathrm{dim}}^3`$, where $`l_{\mathrm{dim}}`$ is typically 100-200. The divergence free condition means we can generate the three magnetic field components for each $`k`$ mode using two complex Gaussian uncorrelated random fields with unit variance, $$𝐂=\left(\begin{array}{c}C_1\\ C_2\end{array}\right).$$ We then determine the magnetic field Fourier components by applying a rotation matrix, $$𝐁=\left(\begin{array}{c}B_x\\ B_y\\ B_z\end{array}\right)=𝐑𝐂,$$ where $`𝐑`$ is a $`3\times 2`$ matrix. From the definition of the magnetic field power spectrum, we see that to get the proper statistical properties, we require $$B_aB_b=R_{am}C_mC_nR_{nb}^T=(𝐑𝐑^T)_{ab}=𝒫(k)P_{ab}(𝐤).$$ While this does not specify the rotation matrix uniquely, it is straightforward to show that choosing the rotation matrix as $$𝐑=\frac{𝒫(k)^{1/2}}{\sqrt{\widehat{k}_x^2+\widehat{k}_y^2}}\left(\begin{array}{cc}\widehat{k}_x\widehat{k}_z& \widehat{k}_y\\ \widehat{k}_y\widehat{k}_z& \widehat{k}_x\\ \left(\widehat{k}_x^2+\widehat{k}_y^2\right)& 0\end{array}\right)$$ (2.9) will produce fields with the correct statistical properties. This rotation is well defined except in the case when $`\widehat{k}_x=\widehat{k}_y=0.`$ Here, $`B_z=0`$ and the other components are uncorrelated, so we instead choose $$𝐑_0=𝒫(k)^{1/2}\left(\begin{array}{cc}1& 0\\ 0& 1\\ 0& 0\end{array}\right).$$ (2.10) The reality of the fields is ensured by requiring $`𝐁(𝐤)=𝐁(𝐤)^{}`$. Throughout, we are careful to avoid creating modes with frequencies higher than the Nyquist frequency of the grid, $`k_{\mathrm{Nyquist}}`$, which could be into power on other frequencies. Since the quantity of greatest interest, the stress-energy, is a quadratic function of the fields, it typically will have power up to twice the cutoff frequency of the magnetic fields. To avoid having aliasing of these fields, we generally require that the magnetic field cutoff frequency be less than half the Nyquist frequency. We also have an infra-red cut-off which is the inevitable result of working on a finite grid. Figure 1 shows a sample Gaussian realization of one component of the magnetic field along a slice through the realization, as well as the resulting trace and traceless components of the stress-energy (the isotropic and anisotropic pressures). Both the isotropic and anisotropic pressures show power on smaller scales than the fields themselves, a direct result of their non-linearity. In addition, the isotropic pressure is darker, reflecting a paucity of positive fluctuations and a significant deviation from Gaussianity. The anisotropic pressure appears to be more similar to the magnetic field – that is, relatively Gaussian. These observations will be made concrete in the next section when we consider the one-point statistics of the pressures. ### II.3 The stress-energy tensor Gravitational waves are sourced directly by the stress-energy of a magnetic field, and the effect from the Lorentz force on the magnetized CMB also depends directly upon this Lewis (2004); for this reason we shall concentrate our study on the space-space part of the magnetic stress-energy tensor, the stresses. In real space, these are $$\tau _{ab}(𝐱)=\frac{1}{2}B_i(𝐱)B_i(𝐱)\delta _{ab}B_a(𝐱)B_b(𝐱).$$ (2.11) In Fourier space, the product of two magnetic field components in real space becomes a convolution, $$\stackrel{~}{\tau }_{ab}(𝐤)=B_a(𝐪)B_b(𝐤𝐪)d^3𝐪.$$ (2.12) The full stress-energy tensor in Fourier space is $$\tau _{ab}(𝐤)=\frac{1}{2}\delta _{ab}\stackrel{~}{\tau }_{ii}(𝐤)\stackrel{~}{\tau }_{ab}(𝐤).$$ (2.13) Even though the stress-energy is non-linear in the fields, we will be assuming throughout that the stress-energy perturbations in the magnetic fields are small compared to the total energy density. Thus, our treatment of the perturbations induced by the magnetic fields in the photons, baryons and so forth are purely linear. In this regime, the evolution of the scalar, vector and tensor perturbations decouple and we will thus decompose the magnetic field sources into these various pieces (interpreted as the isotropic pressure, anisotropic pressure, vortical and anisotropic stresses respectively) as $`\tau _{ab}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\delta _{ab}\tau +(\widehat{k}_a\widehat{k}_b{\displaystyle \frac{1}{3}}\delta _{ab})\tau ^S`$ (2.14) $`+2\widehat{k}_{(a}\tau _{b)}^V+\tau _{ab}^T.`$ We do this in Fourier space by applying combinations of projection operators; $`\tau (𝐤)`$ $`=`$ $`\delta _{ab}\tau _{ab}(𝐤),`$ $`\tau ^S(𝐤)`$ $`=`$ $`\left(\delta _{ab}(3/2)P_{ab}(𝐤)\right)\tau _{ab}(𝐤)Q_{ab}(𝐤)\tau _{ab}(𝐤),`$ (2.15) $`\tau _a^V(𝐤)`$ $`=`$ $`\widehat{k}_{(i}P_{j)a}(𝐤)\tau _{ij}(𝐤)𝒱_{aij}(𝐤)\tau _{ij},`$ $`\tau _{ab}^T(𝐤)`$ $`=`$ $`\left(P_{a(i}(𝐤)P_{j)b}(𝐤){\displaystyle \frac{1}{2}}P_{ij}(𝐤)P_{ab}(𝐤)\right)\tau _{ij}(𝐤)𝒯_{abij}(𝐤)\tau _{ij}(𝐤).`$ In the above, $`(a\mathrm{}b)`$ denotes symmetrization in $`a`$ and $`b`$, i.e, $`A_{(ab)}=(1/2)(A_{ab}+A_{ba})`$, and we are working with symmetrized projectors for the vector and tensor components for future ease. Here, $`\widehat{k}_i\tau _i^V=\widehat{k}_i\tau _{ij}^T=\tau _{ii}^T=0.`$ The first term in the stress-energy contributes purely to the scalar-trace part. For the others we can simply replace $`\tau _{ab}`$ with $`\stackrel{~}{\tau }_{ab}`$. The transfer functions which describe how fluctuations in the magnetic field stress energy impact the microwave background have been previously evaluated, in various semi-analytical limits in Koh and Lee (2000); Durrer et al. (2000); Mack et al. (2002); Subramanian and Barrow (2002) for the temperature perturbations, for example. However, the simplest route towards folding the expected non-Gaussianities onto the CMB, will likely be from the CAMB code as modified recently by Lewis Lewis et al. (2000); Lewis (2004); this produces the transfer functions generated by a magnetic field sourced before neutrino decoupling, neglecting the impact of the scalar perturbations. This could be supplemented by incorporating the extensions to scalar modes of Giovannini Giovannini (2004b) into, for example, the CMBFast code Seljak and Zaldariagga (1996). We leave this issue to a later paper. ## III One-Point Moments There are many ways to characterize non-Gaussianity, particularly given such a strongly non-linear stress-energy term. In this section we briefly consider the skewness and kurtosis of the one-point probability distributions of the isotropic and anisotropic pressures. In this section the results we present are the mean of twenty realizations with a grid-size of $`l_{\mathrm{dim}}=192`$, and the errors quoted are one standard deviation. The simplest to consider is the distribution of the trace part, since it is simply the square of the magnetic field. Despite the divergence-free condition, the three components of the magnetic field at a single point in space are uncorrelated and Gaussian (as was shown above.) Thus we expect the trace of the stress-energy to have a $`\chi ^2`$ distribution with three degrees of freedom. All the moments of a one-point distribution may be given by its moment generating function as $$\mu _n^{}X^n=\frac{^n}{t^n}M(t)\|_{t=0}$$ (3.16) where, for a $`\chi ^2`$ distribution with $`p`$ degrees of freedom $$M(t)=\frac{1}{\left(12t\right)^{p/2}}.$$ (3.17) The central moments are then readily calculated and the normalized skewness and kurtosis are defined to be $$\gamma _1=\frac{\mu _3}{\mu _2^{3/2}},\gamma _2=\frac{\mu _4}{\mu _2^2}3$$ (3.18) We quickly find that, for the $`\chi ^2`$ distribution, the normalized skewness and kurtosis are $$\gamma _1=\sqrt{\frac{8}{p}}1.633,\gamma _2=\frac{12}{p}=4$$ (3.19) where the numerical results are for a distribution with $`3`$ degrees of freedom. The results from the realizations can be seen to be in agreement with the predictions; with a flat spectrum we find that, for the isotropic pressure, $`\gamma _1=1.63\pm 0.01`$ and $`\gamma _2=3.99\pm 0.05`$. For a more realistically observable field, with a power spectrum of $`n=2.5`$, say, we find $`\gamma _1=1.61\pm 0.01`$ and $`\gamma _2=3.92\pm 0.05`$. It is apparent that the statistics for the isotropic pressure are, as expected, relatively insensitive to the spectral index one employs. The anisotropic stress is harder to characterize because it is not a local function of the fields, but contains derivatives of them. However, it effectively is the sum of the products of two Gaussian fields which are, for the most part, independent of each other. The distribution of the product of two independent Gaussians is non-Gaussian but is symmetric (actually following a modified Bessel distribution, as shown in the appendix of Boughn and Crittenden (2005).) Thus the effect of adding such terms is to dilute the skewness. That is, while the isotropic stress is the sum of three very skewed chi-squared distributed variables, the anisotropic stress is the sum of chi-squared terms and symmetric modified Bessels, making the result less skewed. The probability distributions of the isotropic and anisotropic stresses for a flat spectrum are plotted in the left panel of Figure 2 along with a Gaussian and a $`\chi ^2`$ distribution. The damping scale we employed was $`k_c=l_{\mathrm{dim}}/4.1`$i.e. just beneath half the Nyquist frequency. For the steep spectra, we also used twenty realizations and a grid-spacing of $`l_{\mathrm{dim}}=192`$ but with a damping scale at the size of the grid to ensure a reasonable mode coverage in the low-$`k/k_c`$ regime. The anisotropic pressure distribution has quite different properties when the spectral index is changed, including a switch in sign of the skewness. For the flat spectrum we find $`\gamma _1=0.24\pm 0.003`$ and $`\gamma _2=1.10\pm 0.01`$, while with a steep spectrum with a power spectrum of $`n=2.5`$, we find $`\gamma _1=0.38\pm 0.01`$ and $`\gamma _2=0.86\pm 0.02`$. The distributions are plotted in the right-hand side of the figure, again with a sample Gaussian, and the change in the skewness is readily apparent. ## IV Two-Point moments We next calculate the two point power spectra of the various perturbation types. These give a useful example of how the higher order calculations will proceed and give a means of testing our realizations. Some of these have been previously calculated, such as the vector Mack et al. (2002) and tensor Durrer et al. (2000); Mack et al. (2002); Caprini and Durrer (2002) power spectra, while the trace and traceless scalar auto-correlations and their cross correlation, to our knowledge, have not. By the nature of the scalar-vector-tensor decomposition, we do not expect any cross correlations except between the trace and traceless scalar pieces. Thus, we consider five power spectra: one cross spectrum and the auto-spectra of the four pieces of the stress-energy. We focus on constructing rotationally invariant spectra which will contain all the information in the general correlations. In general, the power spectra will involve expectations of four magnetic fields. Since these are assumed to be Gaussian, they can be evaluated using Wick’s theorem which, for four Gaussian fields, may be expressed as $$ABCD=ABCD+ACBD+ADBC.$$ It is most useful to begin with the general two point correlation, $`\stackrel{~}{\tau }_{ab}(𝐤)\stackrel{~}{\tau }_{cd}^{}(𝐩)=\delta (𝐤𝐩){\displaystyle d^3𝐤^{}𝒫(k^{})𝒫(\left\|𝐤𝐤^{}\right\|)}`$ $`\times \left(P_{ac}(𝐤^{})P_{bd}(𝐤𝐤^{})+P_{ad}(𝐤^{})P_{bc}(𝐤𝐤^{})\right).`$ (Note that there are two terms rather than three since we interested in the perturbations from the mean value of the field.) The power spectra of the various stresses may be obtained from this by applying different combinations of the projection operators (2.15) to yield, $$\tau _A(𝐤)\tau _B(𝐩)=\delta (𝐤𝐩)d^3𝐤^{}𝒫(k^{})𝒫(\left\|𝐤𝐤^{}\right\|)_{AB}$$ (4.21) with $`A`$ and $`B`$ denoting the two stress components and $`_{AB}=_{AB}(\gamma ,\mu ,\beta )`$ denoting the relevant angular integrand. The relevant angles possible between the wavevectors have been defined as $$\gamma =\widehat{𝐤}\widehat{𝐤}^{},\mu =\widehat{𝐤}^{}\widehat{𝐤𝐤^{}},\beta =\widehat{𝐤}\widehat{𝐤𝐤^{}},$$ (4.22) where $`\widehat{𝐤𝐤^{}}`$ denotes the unit vector in the direction of $`𝐤𝐤^{}`$. The trace-trace correlation is found by applying the operator $`(1/4)\delta _{ab}\delta _{cd}`$ whence we obtain $$_{\tau \tau }=\frac{1}{2}\left(1+\mu ^2\right).$$ (4.23) Similarly, we obtain the traceless scalar auto-correlation function by applying $`(1)^2Q_{ab}(𝐤)Q_{cd}(𝐤)`$; some algebra yields the result $$F_{\tau ^S\tau ^S}=2+\frac{1}{2}\mu ^2\frac{3}{2}\left(\gamma ^2+\beta ^2\right)3\gamma \mu \beta +\frac{9}{2}\gamma ^2\beta ^2.$$ (4.24) The cross correlation between the trace and traceless scalar pieces requires the operator $`(1/2)\delta _{ab}Q_{cd}(𝐤)`$. This yields $$F_{\tau \tau ^S}=1+\frac{3}{2}\left(\gamma ^2+\beta ^2\right)+\frac{1}{2}\mu ^2\frac{3}{2}\mu \gamma \beta .$$ (4.25) For the vector and tensor contributions, it is useful to construct rotationally invariant combinations that can be relatively easily mapped onto the CMB. The divergenceless condition on the vectors implies that their correlation function can be written as $$\tau _a^V(𝐤)\tau _b^V(𝐩)=\frac{1}{2}𝒫^V(k)P_{ab}(𝐤)\delta (𝐤𝐩).$$ (4.26) where our definition differs by a factor of two from Mack et. al.. All the information is condensed in the rotationally invariant vector isotropic spectrum $`𝒫^V(k)=\tau _a^V(𝐤)\tau _a^V(𝐤).`$ The operator necessary to recover this is $`\widehat{k}_{(a}P_{b)i}(𝐤)\widehat{k}_{(c}P_{d)i}(𝐤)=\widehat{k}_a\widehat{k}_{(c}P_{d)b}(𝐤)+\widehat{k}_d\widehat{k}_{(a}P_{b)c}(𝐤)`$. The resulting angular term can then be shown to be $$F_{\tau ^V\tau ^V}=12\gamma ^2\beta ^2+\mu \gamma \beta .$$ (4.27) Similar arguments apply for the tensor correlations. The full tensor two-point correlation is $$\tau _{ab}^T(𝐤)\tau _{cd}^T(𝐩)=\frac{1}{4}𝒫^T(k)_{abcd}(𝐤)\delta (𝐤𝐩),$$ (4.28) where $`_{abcd}(𝐤)=P_{ac}(𝐤)P_{bd}(𝐤)+P_{ad}(𝐤)P_{bc}(𝐤)P_{ab}(𝐤)P_{cd}(𝐤)`$ which, as can be readily shown, satisfies the transverse-traceless condition on the tensors, and $`\delta _{ac}\delta _{bd}_{abcd}(𝐤)=4`$. We focus on the rotationally invariant tensor isotropic spectrum $`𝒫^T(k)=\tau _{ij}^T(𝐤)\tau _{ij}^T(𝐤)`$. Using the tensor projection operators and simplifying, the relevant operator is $`\left(P_{i(a}(𝐤)P_{b)i}(𝐤)(1/2)P_{ij}(𝐤)P_{ab}(𝐤)\right)\left(P_{i(c}(𝐤)P_{d)i}(𝐤)(1/2)P_{ij}(𝐤)P_{cd}(𝐤)\right)=P_{c(a}(𝐤)P_{b)d}(𝐤)(1/2)P_{ab}(𝐤)P_{cd}(𝐤)`$. This leads to a simple angular term of $$F_{\tau ^T\tau ^T}=(1+\gamma ^2)(1+\beta ^2).$$ (4.29) The vector and tensor results differ from those otherwise presented (see Durrer et. al. and Mack et. al. Durrer et al. (2000); Mack et al. (2002); Caprini and Durrer (2002)), the vector case by $`\beta ^2\gamma ^2`$ as a result of employing the symmetrized projection, and the tensor case by $`\gamma ^2\beta ^2`$. It is straightforward, however, to see that if one redefines the integration wavemode as $`𝐤^{}=𝐤^{\prime \prime }𝐤`$ one maps $`\mu ^{}\mu ,\beta ^{}\gamma ,\gamma ^{}\beta `$. On integration, then, the product $`\gamma ^2\beta ^2`$ is invariant while $`\gamma ^2\beta ^2`$ may be taken to vanish. Our results are thus in agreement with those previously presented. We can compare numerical integrations of these power spectra with the results arising from the realised magnetic fields. Our results for a flat power spectrum ($`n=0`$) are presented in Figure 3 where $`P(k)`$ denotes the various spectra (all presented on the same scale). The agreement between the analytic results (lines) and a simulated field (data points) is striking. We have plotted the power spectra averaged over twenty realizations with a grid-size of $`l_{\mathrm{dim}}=192`$, a damping scale at $`k_c=l_{\mathrm{dim}}/4.1`$, with error bars of 1-$`\sigma `$. We have also rebinned the results because otherwise the data points obscure the theory. We could also consider magnetic fields with $`n>2`$ corresponding to causally-generated fields Caprini and Durrer (2002); the features for such fields are qualitatively similar to those for a flat spectra and there are no difficulties in evaluating the theoretical predictions in this regime. However, as commented, Caprini and Durrer demonstrated that primordial fields of this type would be unobservably small in order to not violate nucleosynthesis bounds on gravitational waves. Moreover, blue spectra naturally tend to pile power around the cut-off scale $`k_c`$ and since we are not at all modelling the microphysics the results would have to be treated with caution. Analytic results can be found in certain limits. There are two regimes of interest for the spectral index, as shown by Durrer et. al. Durrer et al. (2000). For $`n>3/2`$ the integrations are dominated by the cutoff scale, resulting in constant spectra. In this regime, if $`kk_c`$, the angular integrations are straight forward ($`\mu 1,\beta \gamma `$.) Relative to the trace correlation, the amplitudes of the other correlations are $`\tau ^S/\tau =7/5`$, $`\tau _\times /\tau =0`$, $`\tau ^V/\tau =14/15`$ and $`\tau ^T/\tau =28/15`$ respectively. For $`n<3/2`$ the situation is considerably more complex and we content ourselves with the results of our simulations in Figure 4. Immediately apparent is that the effects of the cutoff are reduced by the strongly tilted magnetic spectrum, with each spectrum quickly approaching the power law, $`P_A(k)k^{2n+3}`$, that is naively expected from the $`k`$-integration. Also notable is the change of behaviour of the scalar cross-correlation; whereas this vanishes on large scales in the $`n>3/2`$ regime, it remains finite on large scales for $`n<3/2`$ and so in principle might be observable on the sky. There are also the effects of the infrared cut-off causing a suppression at low-$`k`$. ## V Three-Point Moments In this section we focus on the three point moments in Fourier space, for which it is possible (if laborious) to obtain analytic expressions. The magnetic field realizations provide a way of exploring other kinds of non-Gaussianities which may arise. At the three point level, it is no longer guaranteed that correlations between the scalar, vector and tensor pieces will vanish, and we present some of the first calculations of these here. There are many possible three point moments, but here we consider only the rotationally invariant combinations $`\tau \tau \tau `$, $`\tau \tau \tau ^S`$, $`\tau \tau ^S\tau ^S`$, $`\tau ^S\tau ^S\tau ^S`$, $`\tau \tau _a^V\tau _a^V`$ and $`\tau ^S\tau _a^V\tau _a^V`$. Due to their complexity we postpone calculations of the correlations with the tensor components to a later paper and content ourselves with presenting the results from the realizations; for completeness these are $`\tau \tau _{ab}^T\tau _{ab}^T`$ and $`\tau ^S\tau _{ab}^T\tau _{ab}^T`$ for correlations with the scalars, $`\tau _a^V\tau _{ab}^T\tau _b^V`$ with the vectors and the auto-correlation $`\tau _{ab}^T\tau _{bc}^T\tau _{ac}^T`$. We work throughout in Fourier space, where the three-point moments are known as the bispectra. One advantage of working in Fourier space is that the transfer functions, which fold in the fluid dynamics and describe the impact on the microwave background, are local. ### V.1 General considerations In principle we can calculate all the three point statistics described above in Fourier space. The bispectra involve three wave modes, and since we assume the fields are homogeneous and isotropic, the sum of the three modes must be zero. Thus the bispectra are a function of the amplitudes of the modes alone (or alternatively, two amplitudes and the angle between them.) We denote different geometries by selecting a baseline $`𝐤`$ and a vector $`𝐩`$ making an angle $`\varphi `$ with $`𝐤`$ and having an amplitude $`p=rk`$ (see Figure 5). We may then calculate $`𝐪=𝐤𝐩`$. For simplicity, we here concentrate on the colinear (degenerate) case in which $`𝐩=𝐤`$ implying $`𝐪=2𝐤`$ – that is, $`r=1`$ and $`\varphi =0`$. We calculate the bispectra analogously to the power spectra, although matters are complicated by the need to deal with expectations of six fields rather than four. The object of most general interest is $`_{ijklmn}(𝐤,𝐩,𝐪)\stackrel{~}{\tau }_{ij}(𝐤)\stackrel{~}{\tau }_{kl}(𝐪)\stackrel{~}{\tau }_{mn}(𝐩),`$ which is related to the expectation value of six magnetic fields, $$_{ijklmn}(𝐤,𝐩,𝐪)=d^3𝐤^{}d^3𝐩^{}d^3𝐪^{}B_i(𝐤^{})B_j(𝐤𝐤^{})B_k(𝐩^{})B_l(𝐩𝐩^{})B_m(𝐪^{})B_n(𝐪𝐪^{}).$$ (5.30) As in the two-point case, all three-point moments of interest may be found by applying the relevant projection operator, $`𝒜_{ijklmn}`$, to this. Expanding this six-point correlation with Wick’s theorem generates fifteen terms, eight of which contribute to the reduced bispectrum, that is, the bispectrum neglecting the one-point terms proportional to $`\delta (𝐤)`$, $`\delta (𝐩)`$ or $`\delta (𝐪)`$. This leads eventually to $`_{ijklmn}(𝐤,𝐩,𝐪)=`$ $`\delta (𝐤+𝐩+𝐪){\displaystyle }d^3𝐤^{}𝒫(k^{})𝒫(\|𝐤𝐤^{}\|)𝒫(\|𝐩+𝐤^{}\|)P_{ik}(𝐤^{})\left(P_{jm}(𝐤𝐤^{})P_{ln}(𝐩+𝐤^{})\right)+\{ij,pq\},\{kl,mn.\}`$ These eight terms reduce to the same contribution if the projection tensor $`𝒜_{ijklmn}`$ that recovers a set bispectrum is independently symmetric in $`\{ij\}`$, $`\{kl\}`$ and $`\{mn\}`$. In the power spectra calculations, the geometry was straight forward; here it is considerably more complicated. The three wavevectors of the bispectrum are constrained by homogeneity to obey $`𝐤+𝐪+𝐩=0`$. Combined with the dummy integration wavevector, these define a four sided tetrahedron. This has six edges, $`𝐤,𝐪,𝐩,𝐤^{},𝐤𝐤^{}`$ and $`𝐩+𝐤^{}`$. From these, we can generate fifteen unique angles which could come into the bispectra calculation. This is to be compared to just three edges and three angles required for the power spectra. Clearly these angles are not all independent; they are, in fact, functions of just five underlying angles. For our purposes, it is easiest to work with the fifteen which we separate into four hierarchies; those between the set wavevectors $`𝐤`$, $`𝐩`$ and $`𝐪`$, angle cosines of these vectors with $`𝐤^{}`$, cosines with $`𝐤𝐤^{}`$, and cosines with $`𝐩+𝐤^{}`$. The final group are defined below. We take the angles, with $`\{𝐚,𝐛\}\{𝐤,𝐩,𝐪\}`$, to be $`\theta _{ab}=\widehat{𝐚}\widehat{𝐛},\alpha _a=\widehat{𝐚}\widehat{𝐤}^{},\beta _a=\widehat{𝐚}\widehat{𝐤𝐤^{}},\gamma _a=\widehat{𝐚}\widehat{𝐩+𝐤^{}},`$ $`\overline{\beta }=\widehat{𝐤}^{}\widehat{𝐤𝐤^{}}\overline{\gamma }=\widehat{𝐤}^{}\widehat{𝐩+𝐤^{}}\overline{\mu }=\widehat{𝐤𝐤^{}}\widehat{𝐩+𝐤^{}}.`$ (5.32) In terms of the angles $`\xi _{kq}`$ and $`\xi _{pq}`$ in Figure 5 we obviously have $`\theta _{kq}=\mathrm{cos}\xi _{kq}`$ and similarly for $`\theta _{pq}`$. We also have $`\alpha _k=\mathrm{cos}\overline{\theta }`$. In a manner entirely analogous to the two-point case we find the different bispectra by applying to (V.1) different projection operators to extract the relevant scalar, vector or tensor parts and, given that we ensure that $`𝒜_{ijklmn}`$ has the required symmetries, we may express the bispectra as $$\tau _A(𝐤)\tau _B(𝐪)\tau _C(𝐩)=\delta (𝐤+𝐩+𝐪)d^3𝐤^{}𝒫(k^{})𝒫(\left\|𝐤𝐤^{}\right\|)𝒫(\left\|𝐩+𝐤^{}\right\|)\left(8_{ABC}\right)$$ (5.33) where $`\left\{ABC\right\}`$ denote denote different parts of the stress-energy tensor and $`_{ABC}`$ is the relevant angular component. ### V.2 Scalar bispectra We begin with the simplest case, the bispectrum of the magnetic pressure. This is found by defining $`𝒜_{ijklmn}=(1/8)\delta _{ij}\delta _{kl}\delta _{mn}`$ which gives us $$8_{\tau \tau \tau }=\overline{\beta }^2+\overline{\gamma }^2+\overline{\mu }^2\overline{\beta }\overline{\gamma }\overline{\mu }.$$ (5.34) The first scalar cross-correlation will be between the square of the pressure and the anisotropic pressure, found by using $`𝒜_{ijklmn}=(1/4)\delta _{ij}\delta _{kl}Q_{mn}(𝐪)`$ to give $$8_{\tau \tau \tau ^S}=3\left(1\alpha _q^2\gamma _q^2\beta _q^2\frac{1}{3}(\overline{\beta }^2+\overline{\gamma }^2+\overline{\mu }^2)+\alpha _q(\beta _q\overline{\beta }+\gamma _q\overline{\gamma })+\overline{\mu }(\beta _q\gamma _q+\frac{1}{3}\overline{\beta }\overline{\gamma })\overline{\beta }\overline{\gamma }\beta _q\gamma _q\right)$$ (5.35) Similarly the second scalar cross-correlation, with $`𝒜_{ijklmn}=(1/2)\delta _{ij}Q_{kl}(𝐩)Q_{mn}(𝐪)`$, gives $$_{\tau \tau ^S\tau ^S}=\underset{n=0}{\overset{5}{}}_{\tau \tau ^S\tau ^S}^n$$ (5.36) where $`8_{\tau \tau ^S\tau ^S}^0`$ $`=`$ $`6,`$ $`8_{\tau \tau ^S\tau ^S}^1`$ $`=`$ $`0,`$ $`8_{\tau \tau ^S\tau ^S}^2`$ $`=`$ $`\overline{\beta }^2+\overline{\gamma }^2+\overline{\mu }^2+3\left(\alpha _p^2+\alpha _q^2+\beta _p^2+\beta _q^2+\gamma _p^2+\gamma _q^2\right)+9\theta _{pq}^2,`$ (5.37) $`8_{\tau \tau ^S\tau ^S}^3`$ $`=`$ $`\left(\overline{\beta }\overline{\gamma }\overline{\mu }+3\overline{\mu }(\beta _p\gamma _p+\beta _q\gamma _q)+\overline{\gamma }(\alpha _p\gamma _p+\alpha _q\gamma _q)+\overline{\beta }(\alpha _p\beta _p+\alpha _q\beta _q)+9\theta _{pq}(\alpha _p\alpha _q+\beta _p\beta _q+\gamma _p\gamma _q)\right)`$ $`8_{\tau \tau ^S\tau ^S}^4`$ $`=`$ $`3\left(\overline{\beta }(\overline{\mu }\alpha _p\gamma _p+\overline{\gamma }\beta _q\gamma _q+3\alpha _p\beta _q\theta _{pq})+3\left(\alpha _p\gamma _p\alpha _q\gamma _q+\beta _p\gamma _p\beta _q\gamma _q\right)\right)`$ $`8_{\tau \tau ^S\tau ^S}^5`$ $`=`$ $`9\overline{\beta }\alpha _p\gamma _p\beta _q\gamma _q.`$ Finally the anisotropic scalar bispectrum is found by applying $`A_{ijkmln}=(1)^3Q_{ij}(𝐤)Q_{kl}(𝐩)Q_{mn}(𝐪)`$ which results in $$_{\tau ^S\tau ^S\tau ^S}=\underset{n=0}{\overset{6}{}}_{\tau ^S\tau ^S\tau ^S}^n$$ (5.38) with $`8_{\tau ^S\tau ^S\tau ^S}^0`$ $`=`$ $`9`$ $`8_{\tau ^S\tau ^S\tau ^S}^1`$ $`=`$ $`0`$ $`8_{\tau ^S\tau ^S\tau ^S}^2`$ $`=`$ $`\left(\overline{\beta }^2+\overline{\gamma }^2+\overline{\mu }^2+3(\alpha _k^2+\alpha _p^2+\alpha _q^2+\beta _k^2+\beta _p^2+\beta _q^2+\gamma _k^2+\gamma _p^2+\gamma _q^2)+9(\theta _{kp}^2+\theta _{kq}^2+\theta _{pq}^2)\right)`$ $`8_{\tau ^S\tau ^S\tau ^S}^3`$ $`=`$ $`3(\overline{\mu }(\beta _k\gamma _k+\beta _p\gamma _p+\beta _q\gamma _q+{\displaystyle \frac{1}{3}}\overline{\beta }\overline{\gamma })+\overline{\gamma }(\alpha _k\gamma _k+\alpha _p\gamma _p+\alpha _q\gamma _q)+\overline{\beta }(\alpha _k\beta _k+\alpha _p\beta _p+\alpha _q\beta _q)`$ $`+3\theta _{kp}(\alpha _k\alpha _p+\beta _k\beta _p+\gamma _k\gamma _p)+3\theta _{kq}(\alpha _k\alpha _q+\beta _k\beta _q+\gamma _k\gamma _q)+3\theta _{pq}(\alpha _p\alpha _q+\beta _p\beta _q+\gamma _p\gamma _q)+9\theta _{kp}\theta _{kq}\theta _{pq})`$ $`8_{\tau ^S\tau ^S\tau ^S}^4`$ $`=`$ $`3(\overline{\gamma }\overline{\mu }\alpha _k\beta _k+\overline{\beta }\overline{\mu }\alpha _p\gamma _p+\overline{\beta }\overline{\gamma }\beta _q\gamma _q+3(\alpha _k\beta _k(\alpha _p\beta _p+\alpha _q\beta _q)+\alpha _p\gamma _p(\alpha _k\gamma _k+\alpha _q\gamma _q)+\beta _q\gamma _q(\beta _k\gamma _k+\beta _p\gamma _p))`$ $`+3(\overline{\mu }\theta _{kp}\beta _k\gamma _p+\overline{\gamma }\theta _{kq}\alpha _k\gamma _q+\overline{\beta }\theta _{pq}\alpha _p\beta _q)+9(\theta _{kp}\theta _{kq}\gamma _p\gamma _q+\theta _{kp}\theta _{pq}\beta _k\beta _q+\theta _{kq}\theta _{pq}\alpha _k\alpha _p))`$ $`8_{\tau ^S\tau ^S\tau ^S}^5`$ $`=`$ $`9\left(\overline{\mu }\alpha _k\beta _k\alpha _p\gamma _p+\overline{\gamma }\alpha _k\beta _k\beta _q\gamma _q+\overline{\beta }\alpha _p\gamma _p\beta _q\gamma _q+3(\theta _{kp}\beta _k\gamma _p\beta _q\gamma _q+\theta _{kq}\alpha _k\alpha _p\gamma _p\gamma _q+\theta _{pq}\alpha _k\beta _k\alpha _p\beta _q)\right)`$ $`8_{\tau ^S\tau ^S\tau ^S}^6`$ $`=`$ $`27\alpha _k\beta _k\alpha _p\gamma _p\beta _q\gamma _q.`$ ### V.3 Cross bispectra For the vector and tensor correlations we restrict ourselves to the various rotationally-invariant quantities, which can be identified with cross-correlations between the scalar pressures and either the vector or tensor moduli. The first of these, the correlation between the scalar pressure and the vorticity, we recover with the operator $`𝒜_{ijklmn}=(1/2)\delta _{ij}\widehat{p}_{(k}P_{l)a}(𝐩)\widehat{q}_{(m}P_{n)a}(𝐪)`$. The eventual result is $$_{\tau \tau ^V\tau ^V}=\underset{n=1}{\overset{6}{}}_{\tau \tau ^V\tau ^V}^n$$ (5.39) with $`8_{\tau \tau ^V\tau ^V}^1`$ $`=`$ $`3\theta _{pq}`$ $`8_{\tau \tau ^V\tau ^V}^2`$ $`=`$ $`\gamma _p\gamma _q`$ $`8_{\tau \tau ^V\tau ^V}^3`$ $`=`$ $`\overline{\mu }(\beta _p\gamma _q+\beta _q\gamma _p)+\overline{\gamma }(\alpha _p\gamma _q+\alpha _q\gamma _p)+2\theta _{pq}(\alpha _p^2+\beta _p^2+\gamma _p^2+\alpha _q^2+\beta _q^2+\gamma _q^2+{\displaystyle \frac{1}{2}}\overline{\beta }^2+2\theta _{pq}^2)`$ $`8_{\tau \tau ^V\tau ^V}^4`$ $`=`$ $`\overline{\beta }(\overline{\mu }\alpha _p\gamma _q+\overline{\gamma }\gamma _p\beta _q)2\overline{\beta }\theta _{pq}(\alpha _q\beta _q+\alpha _p\beta _p)2\gamma _p\gamma _q(\alpha _p^2+\beta _p^2+\alpha _q^2+\beta _q^2+{\displaystyle \frac{1}{2}}\overline{\beta }^2+2\theta _{pq}^2)`$ $`2(\alpha _p\alpha _q+\beta _p\beta _q)(\gamma _p^2+\gamma _q^2+2\theta _{pq}^2)`$ $`8_{\tau \tau ^V\tau ^V}^5`$ $`=`$ $`4\theta _{pq}\gamma _p\gamma _q(\alpha _p\alpha _q+\beta _p\beta _q)+2\overline{\beta }\gamma _p\gamma _q(\alpha _p\beta _p+\alpha _q\beta _q)+2\overline{\beta }\alpha _p\beta _q(\gamma _p^2+\gamma _q^2+2\theta _{pq}^2)`$ $`8_{\tau \tau ^V\tau ^V}^6`$ $`=`$ $`4\overline{\beta }\theta _{pq}\alpha _p\gamma _p\beta _q\gamma _q.`$ The cross-correlation with the anisotropic pressure is recovered with the operator $`𝒜_{ijklmn}=Q_{ij}(𝐤)\widehat{p}_{(k}P_{l)a}(𝐩)\widehat{q}_{(m}P_{n)a}(𝐪)`$; the eventual result is $$_{\tau ^S\tau ^V\tau ^V}=\underset{n=1}{\overset{7}{}}_{\tau ^S\tau ^V\tau ^V}^n$$ (5.40) with $`8_{\tau ^S\tau ^V\tau ^V}^1`$ $`=`$ $`6\theta _{pq}`$ $`8_{\tau ^S\tau ^V\tau ^V}^2`$ $`=`$ $`4\gamma _p\gamma _q`$ $`8_{\tau ^S\tau ^V\tau ^V}^3`$ $`=`$ $`\left(\beta _p\gamma _q+\beta _q\gamma _p\right)\overline{\mu }+(\alpha _p\gamma _q+\alpha _q\gamma _p)\overline{\gamma }+\left(3\theta _{kp}\gamma _q+3\theta _{kq}\gamma _p\right)\gamma _k`$ $`+\theta _{pq}\left(6\theta _{kp}^2+6\theta _{kq}^2+4\theta _{pq}^2+3\alpha _k^2+3\beta _k^2+2\alpha _p^2+2\beta _p^2+2\gamma _p^2+2\alpha _q^2+2\beta _q^2+2\gamma _q^2+\beta _b^2\right)`$ $`8_{\tau ^S\tau ^V\tau ^V}^4`$ $`=`$ $`4\theta _{pq}^2\left(3\theta _{kp}\theta _{kq}+\alpha _p\alpha _q+\beta _p\beta _q+\gamma _p\gamma _q\right)+\overline{\beta }\theta _{pq}\left(3\alpha _k\beta _k+2\left(\alpha _p\beta _p+\alpha _q\beta _q\right)\right)`$ $`+6\theta _{pq}\left(\theta _{kp}\left(\alpha _k\alpha _p+\beta _k\beta _p\right)+\theta _{kq}\left(\alpha _k\alpha _q+\beta _k\beta _q\right)\right)`$ $`+\gamma _p\gamma _q\left(3\alpha _k^2+3\beta _k^2+2\alpha _p^2+2\beta _p^2+2\alpha _q^2+2\beta _q+6\theta _{kp}^2+6\theta _{kq}^2+\overline{\beta }^2\right)`$ $`+2\left(\gamma _p^2+\gamma _q^2\right)\left(\alpha _p\alpha _q+\beta _p\beta _q+3\theta _{kp}\theta _{kq}\right)+\overline{\beta }\left(\alpha _p\gamma _q\overline{\mu }+\beta _q\gamma _p\overline{\gamma }\right)+3\left(\alpha _k\gamma _p\overline{\gamma }\theta _{kq}+\beta _k\gamma _q\overline{\mu }\theta _{kp}\right)`$ $`+3\gamma _k\left(\alpha _k\alpha _p\gamma _q+\beta _k\beta _q\gamma _p\right)`$ $`8_{\tau ^S\tau ^V\tau ^V}^5`$ $`=`$ $`4\theta _{pq}^2(\alpha _p\beta _q\overline{\beta }+3\theta _{kq}\alpha _k\alpha _p+3\theta _{kp}\beta _k\beta _q)+4(\theta _{pq}(\alpha _p\alpha _q+\beta _p\beta _q+3\theta _{kp}\theta _{kq})\gamma _p\gamma `$ $`+3\alpha _k\beta _k(\alpha _p\beta _p+\alpha _q\beta _q))+6\theta _{kp}((\alpha _k\alpha _p+\beta _k\beta _p)\gamma _p\gamma _q+\beta _k\beta _q(\gamma _p^2+\gamma _q^2))`$ $`+6\theta _{kq}\left(\left(\alpha _k\alpha _q+\beta _k\beta _q\right)\gamma _p\gamma _q+\alpha _k\alpha _p\left(\gamma _p^2+\gamma _q^2\right)\right)+\left(1\alpha _p\beta _p+1\alpha _q\beta _q+3\alpha _k\beta _k\right)\gamma _p\gamma _q\overline{\beta }`$ $`+2\alpha _p\beta _q\overline{\beta }\left(\gamma _p^2+\gamma _q^2\right)+3\alpha _k\beta _k\left(\alpha _p\gamma _q\overline{\mu }+\beta _q\gamma _p\overline{\gamma }\right)`$ $`8_{\tau ^S\tau ^V\tau ^V}^6`$ $`=`$ $`12\theta _{pq}^2\alpha _k\alpha _p\beta _k\beta _q+4\theta _{pq}\gamma _p\gamma _q\left(\alpha _p\beta _q\overline{\beta }+3\theta _{kq}\alpha _k\alpha _p+3\theta _{kp}\beta _k\beta _q\right)`$ $`+6\alpha _k\beta _k\gamma _p\gamma _q\left(\alpha _p\beta _p+\alpha _q\beta _q\right)+6\alpha _k\beta _k\alpha _p\beta _q\left(\gamma _p^2+\gamma _q^2\right)`$ $`8_{\tau ^S\tau ^V\tau ^V}^7`$ $`=`$ $`12\theta _{pq}\alpha _k\alpha _p\beta _k\beta _q\gamma _p\gamma _q.`$ The full tensor correlation, $`\tau _{ab}^T\tau _{bc}^T\tau _{ac}^T`$, has not been calculated, but it can be found by the application of $`𝒜_{ijklmn}=\left(P_{a(i}(𝐤)P_{j)b}(𝐤)\frac{1}{2}P_{ij}(𝐤)P_{ab}(𝐤)\right)\left(P_{b(k}(𝐩)P_{l)c}(𝐩)\frac{1}{2}P_{kl}(𝐩)P_{bc}(𝐩)\right)\left(P_{a(m}(𝐪)P_{n)c}(𝐪)\frac{1}{2}P_{mn}(𝐪)P_{ac}(𝐪)\right)`$. If we consider first the colinear case, this reduces immediately to a product of projection tensors on $`𝐤`$. Employing the symmetries of $`𝒜_{ijklmn}`$ in $`\{ij\}`$, $`\{kl\}`$ and $`\{mn\}`$ and of $`_{ijklmn}`$ in $`\{ik\}`$, $`\{jm\}`$ and $`\{ln\}`$ one may then demonstrate that this vanishes identically. The general case is not straight forward and we defer it to a later study. ### V.4 Flat spectrum results As was the case for the power spectra, there are two very different spectral regimes for the bispectra. For flat spectra, $`n>1`$, the integrals are dominated by the highest $`k`$ modes around the cutoff scale. For these spectral indices, the bispectra become independent of $`k`$ when $`kk_c`$ and the analytic expressions are straight forward to integrate. Indeed, we can do them exactly in the limit $`kk_c`$; we present these results below both for a generic geometry and for the colinear case. For convenience we define the generic geometry using the two angle cosines $`\alpha _k=\mathrm{cos}(\overline{\theta })`$ and $`\theta _{kq}=\mathrm{cos}(\xi _{kq})`$ from which our specification of $`r`$ and $`\varphi `$ may be recovered – see Figure 5. In the general case, setting $`B=A^2k_c^{3(n+1)}/3(n+1)`$ with $`A`$ the amplitude of the magnetic field power spectrum, we find $`\tau (𝐤)\tau (𝐩)\tau (𝐪)`$ $`=`$ $`B\pi \delta (𝐤+𝐩+𝐪)`$ $`\tau (𝐤)\tau (𝐩)\tau ^S(𝐪)`$ $`=`$ $`0`$ $`\tau (𝐤)\tau ^S(𝐩)\tau ^S(𝐪)`$ $`=`$ $`{\displaystyle \frac{21}{30}}B\pi (2+6\mathrm{cos}(\varphi )\mathrm{cos}(\xi _{kq})3\mathrm{cos}^2(\xi _{kq})3\mathrm{cos}^2(\varphi )+6\mathrm{cos}^2(\varphi )\mathrm{cos}^2(\xi _{kq})`$ $`6\mathrm{cos}(\varphi )\mathrm{cos}^3(\xi _{kq})6\mathrm{cos}^3(\varphi )\mathrm{cos}(\xi _{kq})+6\mathrm{cos}^3(\varphi )\mathrm{cos}^3(\xi _{kq}))\delta (𝐤+𝐩+𝐪)`$ $`\tau ^S(𝐤)\tau ^S(𝐩)\tau ^S(𝐪)`$ $`=`$ $`{\displaystyle \frac{17}{35}}B\pi \left(13\mathrm{cos}^2(\varphi )\mathrm{cos}^2(\xi _{kq})3\mathrm{cos}(\varphi )\mathrm{cos}(\xi _{kq})\mathrm{sin}^2(\varphi )\mathrm{sin}^2(\xi _{kq})\right)\delta (𝐤+𝐩+𝐪)`$ (5.41) $`\tau (𝐤)\tau _a^V(𝐩)\tau _a^V(𝐪)`$ $`=`$ $`{\displaystyle \frac{14}{15}}B\pi (3\mathrm{cos}(\varphi )\mathrm{cos}(\xi _{kq})+\mathrm{sin}^2(\varphi )\mathrm{sin}^2(\xi _{kq})3\mathrm{cos}^3(\varphi )\mathrm{cos}(\xi _{kq})3\mathrm{cos}(\varphi )\mathrm{cos}^3(\xi _{kq})`$ $`+4\mathrm{cos}^3(\varphi )\mathrm{cos}^3(\xi _{kq})\mathrm{sin}^2(\varphi )\mathrm{cos}^2(\varphi )\mathrm{sin}^2(\xi _{kq})\mathrm{sin}^2(\varphi )\mathrm{sin}^2(\xi _{kq})\mathrm{cos}^2(\xi _{kq})`$ $`+4\mathrm{cos}^2(\varphi )\mathrm{cos}^2(\xi _{kq})\mathrm{sin}^2(\varphi )\mathrm{sin}^2(\xi _{kq}))\delta (𝐤+𝐩+𝐪)`$ $`\tau ^S(𝐤)\tau _a^V(𝐩)\tau _a^V(𝐪)`$ $`=`$ $`{\displaystyle \frac{17}{105}}B\pi (6\mathrm{cos}(\varphi )\mathrm{cos}(\xi _{kq})6\mathrm{cos}^3(\varphi )\mathrm{cos}(\xi _{kq})6\mathrm{cos}(\varphi )\mathrm{cos}^3(\xi _{kq})\mathrm{sin}^2(\varphi )\mathrm{sin}^2(\xi _{kq})`$ $`+8\mathrm{cos}^3(\varphi )\mathrm{cos}^3(\xi _{kq})2\mathrm{sin}^2(\varphi )\mathrm{sin}^2(\xi _{kq})\mathrm{cos}^2(\xi _{kq})2\mathrm{sin}^2(\varphi )\mathrm{sin}^2(\xi _{kq})\mathrm{cos}^2(\varphi )`$ $`+8\mathrm{cos}^2(\varphi )\mathrm{cos}^2(\xi _{kq})\mathrm{sin}^2(\varphi )\mathrm{sin}^2(\xi _{kq}))\delta (𝐤+𝐩+𝐪)`$ Specialising these to the colinear case wherein $`\varphi =\xi _{kq}=0`$ these reduce to $`\tau (𝐤)\tau (𝐩)\tau (𝐪)`$ $`=`$ $`B\pi \delta (𝐤+𝐩+𝐪)`$ $`\tau (𝐤)\tau (𝐩)\tau ^S(𝐪)`$ $`=`$ $`0`$ $`\tau (𝐤)\tau ^S(𝐩)\tau ^S(𝐪)`$ $`=`$ $`{\displaystyle \frac{7}{5}}B\pi \delta (𝐤+𝐩+𝐪)`$ $`\tau ^S(𝐤)\tau ^S(𝐩)\tau ^S(𝐪)`$ $`=`$ $`{\displaystyle \frac{34}{35}}B\pi \delta (𝐤+𝐩+𝐪)`$ (5.42) $`\tau (𝐤)\tau _a^V(𝐩)\tau _a^V(𝐪)`$ $`=`$ $`{\displaystyle \frac{14}{15}}B\pi \delta (𝐤+𝐩+𝐪)`$ $`\tau ^S(𝐤)\tau _a^V(𝐩)\tau _a^V(𝐪)`$ $`=`$ $`{\displaystyle \frac{34}{105}}B\pi \delta (𝐤+𝐩+𝐪)`$ The bispectra that are derived from the simulated fields are heavily compromised by the grid-size; unlike the two-point case the three-point moments use only a restricted number of the modes, selected by the geometry chosen for the wavevectors. The result from a single realization is in most cases noise-dominated. To overcome this difficulty we have chosen to simulate a large number of different realizations, taking the mean signal and using their variance to provide an estimate for the errors involved. The results, for a flat power spectrum, a grid-size of $`l_{\mathrm{dim}}=192`$ (and a damping scale of $`k_c=l_{\mathrm{dim}}/4.1`$) and $`1,500`$ combined realizations, are plotted rebinned in Figures 6 with the numerically-integrated predictions overlaid. For simplicity we have concentrated on the colinear case, wherein $`𝐩=𝐤`$ and so $`𝐪=\left(𝐤+𝐩\right)`$. We plot $`k/k_c`$ against $`B(k)`$ where $`B(k)`$ represents the colinear bispectra. ### V.5 Steep spectrum results In the regime $`n<1`$ matters are, as with the two-point case, complicated by the presence of numerous poles; the integrals are dominated by the volume lying between the poles. Rather than attempt a solution, we use our realizations to calculate the bispectra. We present the results for $`n=2.5`$ in Figures 7. There is a dependence on the gridsize for this spectrum due to the paucity of modes of appreciable power given the strongly red spectrum. We tested this by running three realizations with differing grid-sizes, keeping the total number of modes constant in each case constant; specifically we took $`1,500`$ simulations at $`l_{\mathrm{dim}}=192`$, $`5000`$ simulations at $`l_{dim}=128`$, and $`40,000`$ simulations at $`l_{dim}=64`$. As might be expected we found a suppression for the case wherein $`l_{dim}=64`$ – due to scarcity of modes at low $`k`$ – but there was good agreement between the other two cases. We have plotted the results from the $`l_{dim}=192`$ run for greatest dynamic range, again rebinning into $`64`$ bins. With this highly-tilted spectrum we see, as with the two-point case, that the features of the magnetic spectrum at the cut-off scale apparent in the flat case are washed out by the spectral tilt. The predominant shape again is the $`k`$-dependence we expect from the radial component of the integral; in this case $`B(k)k^{3(n+1)}`$. (These are plotted for comparison.) Again, as with the two-point case there is an infrared suppression. The magnitudes of the bispectra fall into three close bands; the strongest are the correlations between the scalars and the tensors, while the middle-band is composed of the scalar auto- and cross-correlations. The weakest bispectra are those involving correlations with the vectors and, in some cases, the $`1\sigma `$ error bars are greater than the mean value; such points have been removed from these plots for aesthetic purposes. The $`\tau _a^V\tau _{ab}^T\tau _b^V`$ correlation is particularly weak. ## VI Discussion and Conclusions We have studied, analytically and via realizations, some of the higher point correlations for tangled large-scale magnetic fields. The analysis is obviously highly model-dependent, and extending the analytic results to other models could be very difficult. However, it would be a simple matter to generalize the numerical realizations to perform the same analysis for a wide variety of time-independent models; all we require is the statistical distribution of the underlying field, the power-spectrum and the form of the stress-energy tensor. For example, while we have assumed the simplest – and perhaps most instructive – case of a magnetic field with Gaussian statistics, it would be a simple matter to employ a $`\chi ^2`$ probability distribution for a magnetic field, which is physically motivated by the creation mechanism considered by Matarrese et. al. Matarrese et al. (2005). (Since our code is time-independent the interpretation would be of the final field immediately prior to recombination.) The realizations will however be limited by the narrow dynamic range allowed in the computation. In our particular case of a tangled primordial magnetic field with a power-law spectrum, we have demonstrated that we can recover the 1-, 2- and 3-point statistics from simulations and with an excellent agreement to our analytic predictions. At the one-point level we have not only verified that the magnetic energy density follows a $`\chi ^2`$ probability distribution function (as expected given that it is directly the square of the underlying Gaussian magnetic field), but that the anisotropic pressure is also non-Gaussian with a significantly more complex relation to the fields. There is also a spectral dependence on its probability distribution function affecting the skewness, which swaps sign as one passes through $`n=3/2`$. The kurtosis, while exhibiting a spectral dependence, remains positive. At the two-point level we have calculated the scalar, vector and tensor auto-correlations, as well as the scalar-cross correlation. For the scale-invariant spectrum we confirm the power-spectra with the expected ratios for scales longer than the cutoff; we also see that the scalar cross-correlation vanishes on large scales but is not in general entirely zero for modes close to the cutoff scale. For a highly-tilted spectrum, the cutoff scale is less important and the spectra behave with power law behavior and with the relative ratios approximately constant. The $`k`$-dependence is the power-law that would be expected from a naive point of view. The surprise is that in this regime the scalar cross-correlation no longer vanishes on large-scales; rather, the correlation remains roughly constant at $`\tau \tau ^S/\sqrt{\tau \tau \tau ^S\tau ^S}0.7.`$ Thus, while the isotropic and anisotropic pressures are indeed correlated, they are not perfectly correlated. At the three-point level we considered a number of rotationally-invariant bispectra, concentrating for simplicity on the colinear case. We find significant non-Gaussianities – in excellent agreement with the analytic predictions. For the flat spectrum, these approach fixed ratios on large scales. These can be both positive or negative, or even zero. As in the two-point case, some qualitative aspects change when we consider a strongly-tilted magnetic power spectrum; features arising from the cutoff scale tend to disappear, leaving instead a simple power-law drop off. The relative ratios of the bispectra also change, and even their relative signs differ. The correlation between the isotropic pressure squared and the anisotropic pressure, which vanishes for scale invariant spectra, becomes non-zero. It remains for future work Brown and Crittenden to fold these results in with the transfer functions for magnetized cosmologies and calculate the non-Gaussianities expected to be imprinted upon the cosmic microwave sky. We can then consider how such signals might be used to constrain the properties of a magnetic field of this type. It would also obviously be straight-forward to consider different magnetic power spectra and statistics. While it is too early to speculate what these will discover, the higher order correlations in the sources will certainly lead to similar higher order correlations in the CMB observables, including perhaps cross correlations between the polarization modes such as $`E^2B^2`$. The techniques used in this paper could be applied to a broad variety of models. Moreover, we have here to our knowledge presented the first calculations of cross-correlations between, for example, scalars and tensors, and demonstrated that they can be of an equal magnitude to scalar auto-correlations. This has great potential relevance to the wider field of sources with non-linear stress-energy tensors, or sources with non-Gaussian initial conditions. For example, it would be interesting to compare our results to the same moments evaluated for defect models, particularly given the current resurgence of interest into networks of cosmic strings. ###### Acknowledgements. We wish to thank A. Lewis, R. Maartens, K. Subramanian and K. Dimopoulos for useful discussions.
warning/0506/physics0506182.html	ar5iv	text	# Determination of Wave Function Functionals: The Constrained-Search—Variational Method ## I Introduction In recent work 1 , we proposed the idea of expanding the space of variations in standard variational calculations of the energy 2 , thereby allowing for an improvement of the energy in such calculations. Equivalently, a required level of accuracy could be achieved with fewer variational parameters. In the traditional application of the variational principle, the space of variations is limited by the choice of analytical form for the approximate wave function. For example, if Gaussian or Slater-type orbitals or a linear combination of such orbitals is employed in the energy functional, the variational space is limited by this choice of functions. The proposed manner by which the space of variations can be expanded is by considering the wave function $`\psi `$ to be a functional of a set of functions $`\chi :\psi =\psi [\chi ]`$, rather than a function. This permits a greater flexibility for the wave function $`\psi [\chi ]`$ because the functions $`\chi `$ may be chosen such that $`\psi [\chi ]`$ reproduces any well-behaved function. In principle, a search over such functions can lead to that function $`\chi `$ for which $`\psi [\chi ]`$ is the true wave function. The space over which the search for the functions $`\chi `$ is to be performed, however, is simply too large for practical purposes, and a subset of this space must be considered. We define the subspace over which the search for the functions $`\chi `$ is to be performed by the requirement that the wave function functional $`\psi [\chi ]`$ satisfy a constraint. Typical constraints on the functional $`\psi [\chi ]`$ are those of normalization, the satisfaction of the Fermi-Coulomb hole sum rule, the requirement that it lead to observables such as the electron density, nuclear magnetic constant, diamagnetic susceptibility, Fermi contact term, or any other physical property of interest. With the wave function functional $`\psi [\chi ]`$ thus determined, a rigorous upper bound to the energy is obtained by application of the variational principle. In this way, not only is a particular property of interest or constarint obtained exactly, the energy is also determined accurately since the variational principle ensures it is correct to second order in the accuracy of the wave function. We refer to this method of determining an approximate wave function as the constrained-search—variational method. The method is general in that it is applicable to both ground and excited states. An attribute of constructing a wave function functional $`\psi [\chi ]`$ via the constrained-search—variational method is that there is an improvement in the structure of the wave function throughout all space. Thus, both single-particle expectations representative of different parts of space as well as two-particle expectations involving two different points in space are obtained accurately. As in standard variational calculations, the satisfaction of constraints imposed on the wave function functional, while ensuring the exactness of a specific property or properties, will nonetheless lead to a less accurate upper bound to the energy provided the space of variations remains fixed. Any such decrease in the accuracy of the upper bound can, however, be offset by an increase in the space of variations. The concept of the wave function $`\psi `$ as a functional $`\psi [\chi ]`$ is general in that the space of variations may be expanded through the functions $`\chi `$. The number of functions $`\chi `$ are also independent of the electron number $`N`$. This contrasts with the Hartree-Fock theory 3 Slater determinant $`\mathrm{\Phi }[\varphi _i]`$ wave function which is also a functional but one of the $`N`$-electron spin-orbitals $`\varphi _i`$. Furthermore, there is no variational-flexibility of these spin-orbitals once they have been determined self-consistently by solution 4 of the Hartree-Fock equations. The space of variations cannot be expanded further, and therefore the Hartree-Fock theory wave function functional cannot be adjusted via the spin-orbitals $`\varphi _i`$ to be the true wave function. Thus, this wave function functional constitutes a point in the variational space as defined for the functional $`\psi [\chi ]`$. The determinantal functional $`\mathrm{\Phi }[\varphi _i]`$ is therefore not general in the manner of the proposed $`\psi [\chi ]`$. In our original work 1 we had noted that the constrained-search—variational method could be extended to the determination of arbitrary Hermitian single-particle operators. In sect.2 we present the equations of this generalization as applied to the ground state of the negative ion of atomic Hydrogen, the Helium atom, and its isoelectronic sequence. The extension of these ideas to excited states in conjunction with the theorem of Theophilou 5 is also described. We had also indicated various ways by which the results presented in our prior work could be improved. One such mechanism was to improve the prefactor in the correlated-determinantal wave function functional. In sect.3 we present the results of the application of the method with such an improved 3-parameter analytical wave function functional to the ground state of the negative ion of atomic Hydrogen and the Helium atom, with normalization as the constraint. We present the results for the total energy $`E`$, the expectations of the Hermitian single-particle operators $`W=_ir_i^n,n=2,1,1,2,W=_i\delta (𝐫_i)`$, and $`W=_i\delta (𝐫_i𝐫)`$, the structure of the dynamic (nonlocal) Coulomb hole charge $`\rho _c(\mathrm{𝐫𝐫}^{})`$ as a function of electron position $`𝐫`$, and the expectations of the two particle operators $`u^2,u,1/u,1/u^2`$, where $`u=\|𝐫_i𝐫_j\|`$. The results for all the expectation values are remarkably accurate when compared with the 1078-parameter wave function of Pekeris 6 , thereby indicating the accuracy of the wave function functionals throughout space. The same accuracy is exhibited in a different way by the comparison of the Coulomb holes with those of the essentially exact holes determined by Slamet and Sahni7 . The results for the energy and two particle expectations are far superior to those of Hartree-Fock theory as expected. However, the single-particle expectations are essentially equivalent since such expectations within Hartree-Fock theory are correct to second order 8 . The comparison with Hartree-Fock theory demonstrates how two square-integrable normalized antisymmetric wave functions can lead to essentially the same electron density9 , but that one can be significantly superior to the other. Our results are also superior to those of the 3-parameter variational Caratzoulas-Knowles wave function 10 that has a similar correlation term as ours but is not a functional. In the concluding section 4, we describe our current work on how the ideas of constructing wave function functionals are being applied in conjunction with Quantal density functional theory11 to the many-electron atom. ## II Constrained- search–variational method In this section we present the generalization of the constrained-search—variational method for constraints whereby typical observables such as the diamagnetic susceptibility, nuclear magnetic constant, Fermi contact term, and the constraint of normalization are determined exactly. For the two-electron systems represented by the negative ion of atomic Hydrogen, the Helium atom, and its isoelectronic sequence, these properties are represented by the expectations of the single-particle operators $`W=r_1^2+r_2^2`$, $`W=1/r_1+1/r_2`$, $`W=\delta (𝐫_1)+\delta (𝐫_\mathrm{𝟐})`$, and $`W=1`$. For these two-electron systems, the Hamiltonian in atomic units ($`e=\mathrm{}=m=1`$) $$\widehat{H}=\frac{1}{2}_1^2\frac{1}{2}_2^2\frac{Z}{r_1}\frac{Z}{r_2}+\frac{1}{r_{12}},$$ (1) where $`𝐫_1`$, $`𝐫_2`$ are the coordinates of the two electrons, $`r_{12}`$ is the distance between them, and $`Z`$ is the atomic number. We next choose the form of the wave function functional to be of the general form $$\psi [\chi ]=\mathrm{\Phi }(s,t,u)[1f(\chi ;s,t,u)],$$ (2) with $`\mathrm{\Phi }(s,t,u)`$ a pre-factor and $`f(\chi ;s,t,u)`$ a correlated correction term: $$f(s,t,u)=e^{qu}(1+qu)[1\chi (q;s,t,u)(1+u/2)],$$ (3) where $`s=r_1+r_2,t=r_1r_2,u=r_{12}`$, are the Hylleraas coordinates12 , and $`q`$ is a variational parameter. Note that any two-electron wave function in a ground or excited state maybe expressed in this form. The key to the wave function functional is the determination of the functions $`\chi (q;s,t,u)`$. The prefactor may be chosen to be of some analytical form with variational parameters as in the present work, or the Hartree-Fock theory wave function 4 , or determined self-consistently within the framework of Quantal Density Functional Theory 11 . For purposes of clarity, and thereby of subsequent analytical ease of solution, we assume the prefactor to depend only on the variables $`s`$ and $`t`$: $`\mathrm{\Phi }=\mathrm{\Phi }(s,t)`$, and for the ground $`1^1S`$ state to be of the analytical form 13 $$\mathrm{\Phi }[\alpha ,\beta ;s,t]=Ne^{\alpha s}cosh(\beta t)=\frac{N}{2}[e^{Z_1r_1}e^{Z_2r_2}+e^{Z_1r_2}e^{Z_2r_1}],$$ (4) where different orbitals are allocated to electrons with up and down spins, $`\alpha `$ and $`\beta `$ are variational parameters, $`Z_1=(\alpha \beta )`$, $`Z_2=(\alpha +\beta )`$, and $`N`$ is the normalization constant (See the Appendix). (Note that the normalization of the prefactor is independent of that of the wave function.) We further assume that $`\chi `$ is a function only of the variable $`s`$: $`\mathrm{\Psi }=\mathrm{\Psi }[\chi (q;s)]`$. (The space of variations could be expanded further by assuming the function $`\chi `$ to depend additionally upon the variable $`t`$, or still further by a dependence on $`t`$ and $`u`$ as well.) The wave function functional $`\mathrm{\Psi }[\chi (q;s)]`$ for the ground state then satisfies the electron-electron cusp condition which in integral form is 14 , $$\mathrm{\Psi }(𝐫_1,𝐫_2,\mathrm{}𝐫_N)=\mathrm{\Psi }(𝐫_2,𝐫_2,𝐫_3,\mathrm{},𝐫_N)(1+r_{12}/2)+𝐫_{12}𝐂(𝐫_2,𝐫_3,\mathrm{},𝐫_N),$$ (5) where $`𝐂(𝐫_2,𝐫_3,\mathrm{},𝐫_N)`$ is an unknown vector. The wave function functional also satisfies the electron-nucleus cusp condition which is 14 , $$\psi (𝐫,𝐫_2,\mathrm{}𝐫_N)=\psi (0,𝐫_2,\mathrm{}𝐫_N).(1Zr)+𝐫𝐚(𝐫_2,\mathrm{}𝐫_N),$$ (6) for $`\alpha =2`$. Here again $`𝐚(𝐫_2,\mathrm{}𝐫_N)`$ is also an unknown vector. In terms of the Hylleraas coordinates, the Hermitian single-particle operators noted above and the normalization operator may be expressed as $`W(s,t)`$ where, respectively, $`W(s,t)=(s^2+t^2)/2`$, $`W(s,t)=\frac{4s}{s^2t^2}`$, $`W(s,t)=\frac{1}{\pi }[\frac{\delta (\frac{(s+t)}{2})}{(s+t)^2}+\frac{\delta (\frac{(st)}{2})}{(st)^2}]`$, and $`W(s,t)=1`$. In general, observables can be represented by single-particle operators expressed as $`W(s,t)`$. The expectation of the operator $`W(s,t)`$ which is $$W=\frac{\mathrm{\Psi }^{}[\chi ]W(s,t)\mathrm{\Psi }[\chi ]𝑑\tau }{\mathrm{\Psi }^{}[\chi ]\mathrm{\Psi }[\chi ]𝑑\tau },$$ (7) can on substitution of the wave function functional $`\psi [\chi ]`$ of Eq.(2) be written as $$\|\mathrm{\Phi }(\alpha ,\beta ;s,t)\|^2[W(s,t)<W>][f^2(q;s,t,u)2f(q;s,t,u)+1]𝑑\tau =0.$$ (8) Equivalently, Eq.(8) may be rewritten as $$_0^{\mathrm{}}e^{2\alpha s}g(s)𝑑s=0,$$ (9) where $$g(s)=_0^s𝑑uu_0^u𝑑tcosh^2(\beta t)[W(s,t)<W>](s^2t^2)[f^2(q;s,t,u)2f(q;s,t,u)+1].$$ (10) We now assume that the expectation $`W`$ is known either through experiment or via some accurate calculation 6 . The next step is the constrained search over functions $`\chi (q;s)`$ for which the expectation $`W`$ of Eq.(7) is obtained. If the parameter $`\alpha `$ in Eq.(9) is fixed, then there exist many functions $`g(s)`$ for which the expectation $`W`$ can be obtained. This corresponds to a large subspace of wave function functionals (See Ref. 1). On the other hand, if the parameter $`\alpha `$ is variable, then the only way in which Eq.(9) can be satisfied is if $$g(s)=0,$$ (11) This is equivalent to the constrained search of all wave function functionals over the subspace in which Eq.(9) is satisfied. Substitution of $`f(\chi ;s,t,u)`$ into Eq.(11) leads to a quadratic equation for the function $`\chi (q;s)`$: $$a(q,s)\chi (q;s)^2+2b(q,s)\chi (q;s)+c(q,s)=0,$$ (12) where $$a(q,s)=_0^s𝑑uu(1+u/2)^2(1+qu)^2e^{2qu}_0^u𝑑tcosh^2(\beta t)(s^2t^2)[W(s,t)<W>],$$ (13) $$b(q,s)=_0^s𝑑u(1+u/2)(1+qu)[e^{2qu}(1+qu)e^{qu}]_0^u𝑑tcosh^2(\beta t)(s^2t^2)[W(s,t)<W>],$$ (14) $$c(q,s)=_0^s𝑑u[e^{2qu}(1+qu)^22e^{qu}(1+qu)+1]_0^u𝑑tcosh^2(\beta t)(s^2t^2)[W(s,t)<W>].$$ (15) Thus, in order to ensure that the wave function functional $`\psi [\chi ]`$ leads to the exact expectation value $`<W(s,t)>`$, one has to solve a quadratic equation for the determination of the functions $`\chi (q;s)`$. The subspace thus corresponds to two points. The two solutions $`\chi _1(q;s)`$ and $`\chi _2(q;s)`$ lead to two normalized wave functions $`\psi [\chi _1]`$ and $`\psi [\chi _2]`$ each of which in turn give rise to the exact expectation $`<W(s,t)>`$. For the two normalized wave function functionals as determined above, the energy functional in terms of Hylleraas coordinates which is $`I[\psi [\chi ]]`$ $`=`$ $`{\displaystyle \psi ^{}\widehat{H}\psi 𝑑\tau }`$ (17) $`=`$ $`2\pi ^2{\displaystyle _0^{\mathrm{}}}ds{\displaystyle _0^s}du{\displaystyle _0^u}dt\{u(s^2t^2)[({\displaystyle \frac{\psi }{s}})^2+({\displaystyle \frac{\psi }{t}})^2+({\displaystyle \frac{\psi }{u}})^2]`$ $`+2{\displaystyle \frac{\psi }{u}}[s(u^2t^2){\displaystyle \frac{\psi }{s}}+t(s^2u^2){\displaystyle \frac{\psi }{t}}]`$ $`[4Zsu(s^2t^2)]\psi ^2\},`$ is then minimized with respect to the parameters $`\alpha `$, $`\beta `$ and $`q`$. The above framework presented for the ground $`1^1S`$ state of the two electron system is general and also applicable to excited states. For example, if one were to consider the excited $`2^3S`$ triplet state of the Helium atom, one could employ for the prefactor in Eq.(2) for the wave function functional $`\psi [\chi ]`$ the expression $`\mathrm{\Phi }(\alpha ;s,t)=\sqrt{\frac{2}{3}}(\frac{\alpha ^4}{\pi })e^{\alpha s}t`$. Note that in this simplest of choices used for explanatory purposes, screening effects are ignored. With such a choice, the procedure to determine the wave function functional $`\psi [\chi ]`$ is the same as described above. In addition, this procedure could be employed in conjunction with the theorem of Theophilou 5 according to which if $`\phi _1,\phi _2,,\phi _m`$,…, are orthonormal trial functions for the $`m`$ lowest eigenstates of the Hamiltonian $`H`$, having exact eigenvalues $`E_1,E_2,E_m`$,… , then $`_{i=1}^m\phi _i\|H\|\phi _i_{i=1}^mE_i`$ . In this way, a rigorous upper bound to the sum of the ground and excited states is achieved. With the ground state energy known, a rigorous upper bound to the excited state energy is then determined, while simultaneously a physical constraint or sum rule is satisfied or an observable obtained exactly. The description of the constrained-search—variational method given in this section concerns the determination of wave function functionals that obtain the expectation value of arbitrary Hermitian single-particle operators exactly. The functions $`\chi `$ were assumed to depend only on the Hylleraas coordinate $`s`$, and as a consequence, a quadratic equation had to be solved for their determination. If the variational space is expanded, then one would have to solve an integral equation for the function $`\chi `$. The ideas of the constrained-search—variational method may also be applied to sum rules involving two-particle properties. For example, consider the pair-correlation density $`g(\mathrm{𝐫𝐫}^{})`$ which is the conditional density at $`𝐫^{}`$ of all other electrons, given that one electron is at $`𝐫`$, and which accounts for electron correlations due to the Pauli exclusion principle and Coulomb repulsion. The pair-correlation density for an N-electron system is defined as $$g(\mathrm{𝐫𝐫}^{})=\mathrm{\Psi }\|\underset{ij}{}\delta (𝐫_i𝐫)\delta (𝐫_j𝐫)\|\mathrm{\Psi }/\rho (𝐫),$$ (18) and satisfies the sum rule $$g(\mathrm{𝐫𝐫}^{})𝑑𝐫^{}=N1,$$ (19) for each electron position $`𝐫`$. However, in order to determine the wave function functional $`\psi [\chi ]`$ that satisfies this sum rule at each electron position, one must solve an integral equation for $`\chi `$. The details of the calculation of such a wave function functional are to be presented elsewhere15 . ## III Application to the ground state of the Helium atom and the negative ion of atomic Hydrogen In this section we apply the constrained-search—variational method as described above to the ground state of the Helium atom and the negative ion of atomic Hydrogen. The constraint employed is that of normalization, and the prefactor is that of Eq. (4). We begin with a discussion of the wave function functionals determined. Wave function functionals The 3-parameter wave function functionals are determined by solution of the quadratic equation Eq.(12). This solution for the functions $`\chi `$ is analytical so that the wave function functionals $`\psi [\chi _1]`$ and $`\psi [\chi _2]`$ too are analytical. We do not provide here the analytical expressions for $`\chi _1(q,\alpha ,\beta ;s)`$ and $`\chi _2(q,\alpha ,\beta ;s)`$, but these functions are plotted in Fig.1. Observe that the two solutions for both $`He`$ and $`H^{}`$ are distinctly different: one is positive and monotonically decreasing while the other is negative and monotonically increasing. Thus, although the two wave functions have the same structural form, and both satisfy the normalization constraint and the electron-electron cusp condition, they are very different. The results as determined by these two wave functions for the ground state energy, and various single- and two-particle expectations are given in the subsections below. Comparisons are made with the results of the prefactor, Hartree-Fock (HF) theory, the 3-parameter Caratzoulas-Knowles (CK), and 1078-parameter Pekeris wave functions. Ground-state energy In Table I, we quote the values for the ground-state energy for $`H^{}`$ and $`He`$. The corresponding satisfaction of the virial theorem, and percent errors when compared to the values of Pekeris for $`He`$ and those of the variational-perturbation results of Aashamar 16 for $`H^{}`$ are also given. Observe that the energies obtained by each wave function functional for $`H^{}`$ and $`He`$ are an order of magnitude superior to that of the prefactor. For $`He`$, these results are on the average $`0.06\%`$ from the Pekeris values. They are also an order of magnitude superior to both those of HF and CK. For $`H^{}`$, both wave function functionals lead to results within $`0.1\%`$ of the Aashamar values, and to positive electron affinities as must be the case since the ion is stable. (In the HF approximation, one does not obtain the negative ion of atomic Hydrogen to be stable. The exact satisfaction of the virial theorem by HF theory, however, is a consequence of self-consistency.) The results clearly demonstrate that highly accurate ground state energies can be obtained by constructing few-parameter wave functions that are functionals. These energies are far superior to those determined by similar wave functions with the same number of parameters but ones that are not functionals. Single-particle expectations In this subsection we present the results of the expectations of the Hermitian single-particle operators $`W=_ir_i^n,n=2,1,1,2,W=_i\delta (𝐫_i)`$, and $`W=_i\delta (𝐫_i𝐫)`$. We begin with the determination of the electron density $`\rho (𝐫)`$, which is the expectation of the operator $`W=_i\delta (𝐫_i𝐫)`$, and from which all the other single-particle expectations may be obtained. (Of course, these expectations may also be determined directly from the wave function functionals.) The density $`\rho (𝐫)`$ is also required for the determination of the nonlocal Coulomb hole charge distribution $`\rho _c(\mathrm{𝐫𝐫}^{})`$ as explained in the following subsection. Now the wave function functionals are in terms of the Hylleraas coordinates $`(s,t,u)`$ which involve the position of both the electrons or both their radial distances from the nucleus. The electron density $`\rho (𝐫)`$, on the other hand, depends only on the coordinates of one of the particles. Its determination from wave functions that are written in terms of the Hylleraas coordinates is as follows. The electron density $$\rho (𝐫)=\psi ^{}(\underset{i}{}\delta (𝐫_i𝐫))\psi 𝑑\tau =2\psi ^2(𝐫𝐫^{})𝑑𝐫^{},$$ (20) Using the symmetry of the two electronic system, we have $$𝑑𝐫^{}=2\pi _0^{\mathrm{}}r^2𝑑r^{}_1^1𝑑cos\theta .$$ (21) With $`u=\sqrt{r^2+r^22rr^{}cos\theta }`$, then, for fixed $`r`$ and $`r^{}`$, we can rewrite Eq.(21) as $$𝑑𝐫^{}=2\pi _0^{\mathrm{}}\frac{r^{}}{r}𝑑r^{}_{\|rr^{}\|}^{r+r^{}}u𝑑u.$$ (22) On rewriting the wave function in terms of ($`r,r^{},u`$), and substituting Eq.(22) into Eq.(20) leads to $`\rho (𝐫)`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{r^{}}{r}}𝑑r^{}{\displaystyle _{\|rr^{}\|}^{r+r^{}}}u\psi ^2(r,r^{},u)𝑑u`$ , $`=`$ $`\rho _0(𝐫)+\mathrm{\Delta }\rho _0(𝐫),`$ (23) where $`\rho _0(𝐫)`$ is the density due to the prefactor (see the Appendix for the analytical expression): $`\rho _0(𝐫)`$ $`=`$ $`2N^2{\displaystyle e^{2\alpha s}cosh^2(\beta t)𝑑𝐫^{}},`$ (24) and $`\mathrm{\Delta }\rho _0(𝐫)`$ $`=`$ $`2N^2{\displaystyle e^{2\alpha s}cosh^2(\beta t)(f^2(x;s,t,u)2f(\chi ;s,t,u))𝑑𝐫^{}},`$ is the density due to the correlation term, which can be evaluated numerically. The electron density at the nucleus is $`\rho (0)`$ $`=`$ $`{\displaystyle }\psi ^{}({\displaystyle \underset{i}{}}\delta (𝐫_i)\psi d\tau `$ (26) $`=`$ $`\rho _0(0)+\mathrm{\Delta }\rho _0(0),`$ where $`\rho _0(0)`$ is the prefactor contribution(see Appendix): $$\rho _0(𝐫)=2N^2e^{2\alpha r}cosh^2(\beta r)𝑑𝐫,$$ (27) and the correlation contribution is $$\mathrm{\Delta }\rho _0(𝐫)=2N^2e^{2\alpha r^{}}cosh^2(\beta r^{})(f^2(x;s,t,u)2f(\chi ;s,t,u))\|_{r_1=r=0,u=r_2=r^{}}d𝐫^{}.$$ (28) In Table II we quote the expectations of the operators $`W=_ir_i^n,n=2,1,1,2`$, and $`W=_i\delta (𝐫_i)`$, for the ground state of the He atom as determined by the functionals $`\psi [\chi _1]`$ and $`\psi [\chi _2]`$ together with those of Hartree-Fock theory, and the Caratzoulas-Knowles and Pekeris wave functions. The corresponding percent errors relative to the values of Pekeris are given in Table III. As expected (see Table III), the improvement over the prefactor values is significant. The results of the two wave function functionals and those of Hartree-Fock theory are essentially equivalent, indicating thereby that the corresponding densities are also essentially the same. The expectations of single-particle operators in Hartree-Fock theory are, of course, known to be correct to second order8 . Hence, both the wave function functionals are accurate throughout space including the deep interior and far exterior of the atom. The comparison with the Caratzoulas-Knowles values (see Table III) is interesting for its implications. The wave function functional values are an order of magnitude superior. Of course, one does not expect the CK results to be accurate because these single-particle expectations are correct only to first order in the accuracy of the wave function. Thus, our results once again demonstrate, that wave function functionals determined by the constrained-search— variational method are superior to variationally determined wave functions that are not functionals. Structure of Coulomb holes We next consider the structure of the Coulomb hole charge distribution $`\rho _c(\mathrm{𝐫𝐫}^{})`$ as a function of the electron position $`𝐫`$. The definition of this nonlocal or dynamic charge whose structure changes with electron position for nonuniform electron gas systems derives from that of the pair-correlation density $`g(\mathrm{𝐫𝐫}^{})`$ of Eq.(18) and from local effective potential energy theory 11 . The pair-density may be separated into its local and nonlocal components as $$g(\mathrm{𝐫𝐫}^{})=\rho (𝐫^{})+\rho _{xc}(\mathrm{𝐫𝐫}^{}),$$ (29) where $`\rho _{xc}(\mathrm{𝐫𝐫}^{})`$ is the Fermi-Coulomb hole charge. This dynamic charge distribution is the change in the pair density relative to the density that occurs as a consequence of the Pauli exclusion principle and Coulomb repulsion. It follows from Eq.(19) that its total charge is $`1`$. The definition of the Coulomb hole $`\rho _{xc}(\mathrm{𝐫𝐫}^{})`$ derives in turn from that of the Fermi-Coulomb $`\rho _{xc}(\mathrm{𝐫𝐫}^{})`$ and Fermi $`\rho _x(\mathrm{𝐫𝐫}^{})`$ holes, the latter being defined through local effective potential energy theory. In this theory, the interacting system as described by the Schrödinger equation is replaced by one of noninteracting Fermions with the same density. The corresponding wave function is a Slater determinant of single-particle spin orbitals, and one can then write down the resulting pair-correlation density $`g_s(\mathrm{𝐫𝐫}^{})`$ of the model system as $$g_s(\mathrm{𝐫𝐫}^{})=\rho (𝐫^{})+\rho _x(\mathrm{𝐫𝐫}^{}),$$ (30) where $`\rho _x(\mathrm{𝐫𝐫}^{})`$, the Fermi hole, is the nonlocal component of this pair density, and is a consequence solely of the Pauli principle. The total charge of the Fermi hole is also $`1`$. The Coulomb hole is then defined as the difference between the Fermi-Coulomb and Coulomb holes: $$\rho _c(\mathrm{𝐫𝐫}^{})=\rho _{xc}(\mathrm{𝐫𝐫}^{})\rho _x(\mathrm{𝐫𝐫}^{}),$$ (31) and is thus representative solely of Coulomb correlations. The total charge of the Coulomb hole is $`0`$. For two-electron systems in local effective potential theory 11 , the Fermi hole is then $`\rho _x(\mathrm{𝐫𝐫}^{})=\rho (𝐫^{})/2`$ independent of electron position $`𝐫`$. In Figs. 2-4, we plot cross sections of the Coulomb hole $`\rho _c(\mathrm{𝐫𝐫}^{})`$ for different electron positions $`𝐫`$ as obtained via the functional $`\psi [\chi _2]`$ together with the ‘exact’ Coulomb hole determined by Slamet and Sahni7 . (The electron, indicated by the arrow, is on the z axis corresponding to $`\theta =0^0`$. The cross section through the Coulomb hole plotted corresponds to $`\theta ^{}=0^0`$ with respect to the electron-nucleus direction. The graph for $`r^{}<0`$ corresponds to the structure for $`\theta ^{}=\pi `$ and $`r^{}>0`$.) The electron positions are at $`r=0,0.566,0.8,1.0,1.5`$, and $`5.0`$ (a.u.). It is evident from these figures that the Coulomb holes as determined from the functional $`\psi [\chi _2]`$ closely approximate the exact results for electron positions throughout space: in the interior, within the atom, near its surface and outside the atom, and in the far asymptotic region. Note the cusp representative of the electron-electron cusp condition at the electron position which is indicated by an arrow in the figures. Two-particle expectations As a consequence of the accuracy of the dynamic Coulomb holes obtained, we expect the results for the expectation of two-particle operators to also be accurate. In Table IV we quote the values for the expectations of the operators $`u^2,u,1/u,1/u^2`$, where $`u=\|𝐫_i𝐫_j\|`$, together with the Hartree-Fock and Pekeris values. The corresponding percent errors compared to those of Pekeris are given in Table V. Once again, the results are an order of magnitude superior to those of the prefactor, and are accurate for both functionals, although those due to $`\psi [\chi _2]`$ are consistently superior (see Table V). Of course, as expected, the Hartree-Fock theory results are not accurate. If one were able to write the expectation of arbitrary operators $`\widehat{O}`$ as functionals of the density: $`\widehat{O}=\psi [\rho ]\|\widehat{O}\|\psi [\rho ]=O[\rho ]`$, as is possible in principle according to the Hohenberg-Kohn theorem 17 , then it is in the expectation of two-particle operators that the small differences between the Hartree-Fock theory density and those of the two wave function functionals would be exhibited. ## IV Concluding remarks The idea of expanding the space of variations in variational calculations by writing the wave function as a functional of functions is appealing not only because the functionals lead to more accurate upper bounds for the energy with fewer parameters, but also because, as demonstrated in this work, they lead to wave functions that are accurate over all space. Thus, both single- and two-particle expectations are also determined accurately. Certainly, one could claim by comparison with the results of Hartree-Fock theory, but without rigorous proof, that single-particle expectations obtained thereby are correct to second order in the accuracy of the wave function. It is also evident that the accuracy of two-particle expectations lies somewhere between first and second order. In contrast, variationally determined wave functions that are not functionals are accurate only in those regions of space contributing to the energy. Thus, for such wave functions, it is the expectation value of only those single- and two-particle operators that appear in the Hamiltonian that are reasonably accurate. All other expectations are correct only to first order. The results of the present work could be further improved as follows: by expanding the space of variations through the function $`\chi `$; by employing other more efficacious choices for the analytical form of the correlation factor and thus of the wave function functional; and by improving the prefactor. In our work so far, we have employed analytical forms for the prefactor. ( The results of our prefactor for the ground state energy of both $`H^{}`$ and $`He`$ are superior to those of Hartree-Fock theory, see Table I.) Of course, one could employ the Hartree-Fock theory Slater determinant as the prefactor. Or one could employ a determinantal prefactor based on the orbitals generated within the local effective potential framework of Quantal density functional theory (Q-DFT). In principle, these orbitals generate the true electron density via a model system of noninteracting Fermions. The corresponding local potential within Q-DFT depends upon the wave functions of the interacting and noninteracting systems. Therefore, the corresponding orbitals generated are representative of electron correlations due to the Pauli exclusion principle, Coulomb repulsion, and the correlation contributions to the kinetic energy. Finally, we are presently investigating the use of wave function functionals in conjunction with Q-DFT for the many-electron case of $`N>2`$. In these calculations, the antisymmetric determinantal correlated wave function functional employed is of the form $$\psi [\chi ]=\mathrm{\Phi }\{\varphi _i\}\mathrm{\Pi }_{ij}(1f(\chi ;𝐫_i,𝐫_j)).$$ (32) Here $`\mathrm{\Phi }\{\varphi _i\}`$ is a Slater determinant that defines the state of the system and whose orbitals $`\varphi _i`$ are generated via the differential equation of Q-DFT, $`f(\chi ;𝐫_i,𝐫_j)`$ is a spinless correlation functional: $`f(\chi ;𝐫_i,𝐫_j)=e^{\beta ^2r^2}[1\chi (R)(1+r/2)]`$, where $`𝐫=𝐫_i𝐫_j,𝐑=𝐫_i+𝐫_j,\beta =q\rho ^{1/3}(R)`$, $`q`$ is a variational parameter, and $`\chi (R)`$ is determined by the constraint of the Coulomb hole sum rule for each electron position. This wave function functional satisfies the electron-electron cusp condition. In this instance an integral equation is solved 15 to determine the function $`\chi (R)`$. Further, the products of the correlation functional are limited to lowest order since higher order products of these factors are less significant 18 . The highest occupied eigenvalue of Q-DFT differential equation corresponds in principle to the negative of the ionization potential11 . The region that contributes principally to this eigenvalue is the asymptotic classically forbidden region of the atom. In Q-DFT, the asymptotic structure of the effective potential is due solely to Pauli correlations, and can be determined exactly. This is because the contributions to the potential due to Coulomb correlations and Correlation-Kinetic effects decay more rapidly than $`(1/r)`$11 , so that the potential in this region arises only from the Fermi hole charge which is defined through the Slater determinant of the orbitals. Thus, accurate ionization potentials cab be obtained via the use of correlated-determinantal wave function functionals in conjunction with Q-DFT. These are variational-self—consistent calculations that lead to upper bounds for the energy while simultaneously satisfying a nonlocal physical constraint. We are also currently investigating the construction of wave function functionals of the form employed in the present work, but with the satisfaction of constraints other than that of normalization. * ## Appendix A We give the analytical expressions for the normalization constant, the energy, and various single- and two-particle expectation values as determined by the prefactor wave function $$\mathrm{\Phi }=Ne^{\alpha s}cosh(\beta t).$$ (33) Normalization $`{\displaystyle 𝑑\tau \mathrm{\Phi }^2}`$ $`=`$ $`2\pi ^2N^2{\displaystyle _0^{\mathrm{}}}𝑑se^{2\alpha s}{\displaystyle _0^s}𝑑tcosh^2(\beta t){\displaystyle _t^s}𝑑uu(s^2t^2)`$ (34) $`=`$ $`N^2\pi ^2({\displaystyle \frac{2\alpha ^6+3\alpha ^4\beta ^23\alpha ^2\beta ^4+\beta ^6}{2\alpha ^6(\beta \alpha )^3(\alpha +\beta )^3}})=1.`$ Ground-state energy $`E_0`$ $`=`$ $`{\displaystyle \mathrm{\Phi }^{}\widehat{H}\mathrm{\Phi }𝑑\tau }`$ (35) $`=`$ $`2\pi ^2{\displaystyle _0^{\mathrm{}}}ds{\displaystyle _0^s}du{\displaystyle _0^u}dt\{u(s^2t^2)[({\displaystyle \frac{\mathrm{\Phi }}{s}})^2+({\displaystyle \frac{\mathrm{\Phi }}{t}})^2+({\displaystyle \frac{\mathrm{\Phi }}{u}})^2]`$ $`+2{\displaystyle \frac{\mathrm{\Phi }}{u}}[s(u^2t^2){\displaystyle \frac{\mathrm{\Phi }}{s}}+t(s^2u^2){\displaystyle \frac{\mathrm{\Phi }}{t}}][4Zsu(s^2t^2)]\mathrm{\Phi }^2\}`$ $`=`$ $`\alpha ^22Z\alpha +{\displaystyle \frac{\alpha (\beta ^2\alpha ^2)(10\alpha ^411\alpha ^2\beta ^2+5\beta ^4)}{8(2\alpha ^6+3\alpha ^4\beta ^23\alpha ^2\beta ^4+\beta ^6)}}`$ $`{\displaystyle \frac{\beta ^4(3\alpha ^43\alpha ^2\beta ^2+\beta ^4)}{(2\alpha ^6+3\alpha ^4\beta ^23\alpha ^2\beta ^4+\beta ^6)}}.`$ Expectation values $$\rho _0(𝐫)=\delta (𝐫_1𝐫)+\delta (𝐫_2𝐫)=N^2\pi e^{2\alpha r_1}[\frac{1}{\alpha ^3}+\frac{1}{2}e^{2\beta r_1}(\frac{1}{(\alpha \beta )^3}+\frac{e^{4\beta r_1}}{(\alpha +\beta )^3})],$$ (36) $$\rho _0(0)=\delta (𝐫_1)+\delta (𝐫_2)=N^2\pi [\frac{1}{\alpha ^3}+\frac{1}{2}(\frac{1}{(\alpha \beta )^3}+\frac{1}{(\alpha +\beta )^3})],$$ (37) $`r_1+r_2={\displaystyle 𝑑\tau s\mathrm{\Phi }^2}`$ $`=`$ $`2\pi ^2N^2{\displaystyle _0^{\mathrm{}}}𝑑sse^{2\alpha s}{\displaystyle _0^s}𝑑tcosh^2(\beta t){\displaystyle _t^s}𝑑uu(s^2t^2)`$ (38) $`=`$ $`N^2\pi ^2({\displaystyle \frac{3(2\alpha ^84\alpha ^6\beta ^2+6\alpha ^4\beta ^44\alpha ^2\beta ^6+\beta ^8)}{2\alpha ^7(\beta \alpha )^4(\alpha +\beta )^4}}).`$ $`{\displaystyle \frac{1}{r_1}}+{\displaystyle \frac{1}{r_2}}={\displaystyle 𝑑\tau \frac{4s}{s^2t^2}\mathrm{\Phi }^2}`$ $`=`$ $`2\pi ^2N^2{\displaystyle _0^{\mathrm{}}}𝑑s4se^{2\alpha s}{\displaystyle _0^s}𝑑tcosh^2(\beta t){\displaystyle _t^s}𝑑uu`$ (39) $`=`$ $`2\alpha .`$ $`{\displaystyle \frac{1}{r_1^2}}+{\displaystyle \frac{1}{r_2^2}}={\displaystyle 𝑑\tau \frac{8(s^2+t^2)}{(s^2t^2)^2}\mathrm{\Phi }^2}`$ $`=`$ $`2\pi ^2N^2{\displaystyle _0^{\mathrm{}}}𝑑se^{2\alpha s}{\displaystyle _0^s}𝑑tcosh^2(\beta t)4(s^2+t^2)`$ (40) $`=`$ $`N^2\pi ^2{\displaystyle \frac{(4\alpha ^64\alpha ^4\beta ^2+6\alpha ^2\beta ^42\beta ^6)}{\alpha ^4(\beta \alpha )^3(\alpha +\beta )^3}}.`$ $`r_1^2+r_2^2`$ $`=`$ $`{\displaystyle 𝑑\tau \frac{(s^2+t^2)}{2}\mathrm{\Phi }^2}=N^2\pi ^2{\displaystyle _0^{\mathrm{}}}𝑑se^{2\alpha s}{\displaystyle _0^s}𝑑t{\displaystyle \frac{cosh^2(\beta t)(s^2+t^2)(s^2t^2)^2}{2}}`$ (41) $`=`$ $`N^2\pi ^2{\displaystyle \frac{3(2\alpha ^{10}+4\alpha ^8\beta ^210\alpha ^6\beta ^4+10\alpha ^4\beta ^65\alpha ^2\beta ^8+\beta ^{10})}{\alpha ^8(\beta \alpha )^5(\alpha +\beta )^5}}.`$ $`r_{12}={\displaystyle 𝑑\tau u\mathrm{\Phi }^2}`$ $`=`$ $`2N^2\pi ^2{\displaystyle _0^{\mathrm{}}}𝑑se^{2\alpha s}{\displaystyle _0^s}𝑑tcosh^2(\beta t)(s^2t^2){\displaystyle _t^s}u^2𝑑u`$ (42) $`=`$ $`{\displaystyle \frac{(70\alpha ^8126\alpha ^6\beta ^2+209\alpha ^4\beta ^4140\alpha ^2\beta ^6+35\beta ^8)}{16\alpha (\beta ^2\alpha ^2)(2\alpha ^6+3\alpha ^4\beta ^23\alpha ^2\beta ^4+\beta ^6)}}.`$ $`r_{12}^2={\displaystyle 𝑑\tau u^2\mathrm{\Phi }^2}`$ $`=`$ $`2N^2\pi ^2{\displaystyle _0^{\mathrm{}}}𝑑se^{2\alpha s}{\displaystyle _0^s}𝑑tcosh^2(\beta t)(s^2t^2){\displaystyle _t^s}u^3𝑑u`$ (43) $`=`$ $`{\displaystyle \frac{6(2\alpha ^{10}+4\alpha ^8\beta ^210\alpha ^6\beta ^4+10\alpha ^4\beta ^65\alpha ^2\beta ^8+\beta ^{10})}{\alpha ^2(\beta ^2\alpha ^2)^2(2\alpha ^6+3\alpha ^4\beta ^23\alpha ^2\beta ^4+\beta ^6)}}.`$ $`{\displaystyle \frac{1}{r_{12}}}={\displaystyle 𝑑\tau \frac{1}{u}\mathrm{\Phi }^2}`$ $`=`$ $`2N^2\pi ^2{\displaystyle _0^{\mathrm{}}}𝑑se^{2\alpha s}{\displaystyle _0^s}𝑑tcosh^2(\beta t)(s^2t^2){\displaystyle _t^s}𝑑u`$ (44) $`=`$ $`{\displaystyle \frac{\alpha (\beta ^2\alpha ^2)(10\alpha ^411\alpha ^2\beta ^2+5\beta ^4)}{8(2\alpha ^6+3\alpha ^4\beta ^23\alpha ^2\beta ^4+\beta ^6)}}.`$ $$\frac{1}{r_{12}^2}=𝑑\tau \frac{1}{u^2}\mathrm{\Phi }^2=2N^2\pi ^2_0^{\mathrm{}}𝑑se^{2\alpha s}_0^s𝑑tcosh^2(\beta t)(s^2t^2)_t^s\frac{1}{u}𝑑u$$ (45) Eq.(A.13) can be evaluated numerically. ###### Acknowledgements. This work was supported by the Research Foundation of CUNY. L. M. was supported in part by NSF through CREST, and by a “Research Centers in Minority Institutions” award, RR-03037, from the National Center for Research Resources, National Institutes of Health.
warning/0506/astro-ph0506356.html	ar5iv	text	# VLT-SINFONI observations of Mrk 609 - A showcase for X-ray active galaxies chosen from a sample of AGN suitable for adaptive optics observations with natural guide stars ## 1 Introduction A major cornerstone for extragalactic astronomy is the advent of adaptive optics (AO) assisted imaging and spectroscopy on large ground-based telescopes like the Very Large Telescope (VLT), offering a combination of (near) diffraction-limited resolving power and large light-collecting area of 8-10m class telescopes (e.g. Brandner & Kasper, 2005). In addition 3D spectroscopy allows to study both, the morphology and the chemical composition, as well as the dynamics of extragalactic objects at the same time, at an unprecedented depth. ### 1.1 Starburst/Seyfert composite galaxies Despite the finding of an apparent coevolution of super massive black holes in the centers of galaxies and their galaxy bulges (hosts) (e.g. Page et al., 2001) the detailed nature of this interconnection remains mysterious. Is star formation triggered by the active galactic nucleus (AGN) due to radiation pressure and winds from the accretion disk, which disturb the interstellar medium (as discussed by van Breugel & Dey (1993) for 3C 285)? Or can a nuclear starburst component initiate the accretion process onto the black hole (Norman & Scoville, 1988)? Recently, in their classification study of IRAS selected ROSAT sources, Moran et al. (1996) discovered a class of starburst/Seyfert composite galaxies. They show optical spectra which are dominated by features of starburst galaxies, using the line diagnostics of Veilleux & Osterbrock (1987). Their X-ray luminosity, however, is typical for Seyfert 2 galaxies. A closer look at the spectra reveals some Seyfert-like features, e.g. \[OIII\]$`\lambda 5007`$ is significantly broader than all the narrow emission lines in the optical spectrum and in some cases there is a weak broad H$`\alpha `$ component. There appears a resemblance with narrow-line X-ray galaxies (e.g. Boyle et al., 1995), which also show spectra of composite nature. It is still not clear how their strong X-ray emission can be reconciled with the weak/absent optical Seyfert characteristics. The faintness of these objects in the X-ray as well as the optical domain did not allow to study them in detail so far. Near infrared (NIR) studies, especially integral field spectroscopy, provide powerful means to investigate the (circum-) nuclear properties of the above described AGN. Besides the much lower dust extinction there are a number of NIR diagnostic lines (in emission as well as in absorption) to probe the excitation mechanisms and stellar populations in these objects. Among these are hydrogen recombination lines, rotational/vibrational transitions of H<sub>2</sub>, stellar features like the CO(2-0) and CO(6-3) absorption band heads and forbidden lines like \[FeII\] and \[SiVI\] (Hill et al., 1999; Mouri, 1994; Marconi et al., 1994). ### 1.2 Mining the sky: A sample of X-ray bright AGN Multi wavelength sky surveys like the Sloan Digital Sky Survey (SDSS) and the ROSAT All Sky Survey (RASS) are very comprehensive databases to search in for targets suitable for very sensitive and high-resolution AO-assisted observations in the NIR. In particular, a cross-correlation between the first data release of the SDSS (Abazajian et al., 2003) and the RASS (Voges et al., 1999) resulted in a sample of about 70 X-ray luminous AGN ($`L_X10^{43}10^{45}`$ erg s<sup>-1</sup>) at redshifts between $`z=0.1`$ and $`z=1`$ (Zuther et al., 2004, 2005). Most optical counterparts of the X-ray sources turn out to be AGN (e.g. Giacconi et al., 2001). Furthermore, these X-ray luminous sources cannot be studied locally because of their small number density in the local universe (Hasinger, 1998). They are therefore ideal targets for adaptive optics observations. A subset of this sample is comprised of the composite galaxies described in Section 1.1. Our sample allows to study their optical/near-infrared properties in the above described framework. In the following we will present our first integral field observations of the composite galaxy Mrk 609 (Fig. 1). ## 2 SINFONI observations of Mrk 609 From the subset of starburst/Seyfert composite galaxies we chose Mrk 609 (Rudy et al., 1988) (Fig. 1) as one of the closest and brightest (AO self referencing) objects for one of the science verification phase observations<sup>1</sup><sup>1</sup>1http://www.eso.org/science/vltsv/sinfonisv/xrayagn.html of the new AO-assisted integral field spectrometer SINFONI at the Very Large Telescope (Eisenhauer et al., 2003). The observations presented in this contribution were taken in AO-mode with a 100 mas pixel scale and a field of view of 3$`\times `$3 arcsec<sup>2</sup>. The 2-dimensional image was sliced by small mirrors into 32 slitlets which then were reimaged onto one long pseudoslit and dispersed onto a 2k$`\times `$2k detector. Using the filters $`J`$ and $`H+K`$, a spectral resolution of $`R2000`$ was achieved. Details of the data reduction will be presented in a forthcoming paper (Zuther et al. in prep.). In the following we will present very first results from this study. ### 2.1 Overall properties This dataset shows the wealth of information which can be retrieved from the integral field observation. Fig. 2 shows a median $`H+K`$ continuum image of the reconstructed 3-dimensional data cube. Shown are the inner $`3^{\prime \prime }\times 3^{\prime \prime }`$ around the nucleus. At its redshift of $`z=0.034`$, 1 arcsecond corresponds to about 700 pc<sup>2</sup><sup>2</sup>2Assuming H$`{}_{0}{}^{}=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, and $`\mathrm{\Omega }_m=1`$. and 1 pixel ($``$ 0.05 arcsec) to about 40 pc. The shape of the contours, which appear to be elongated towards the root points of the spiral arms (Fig. 1), suggests the possibility of the presence of a nuclear bar (cf. e.g. Martini et al, 2001). Fig. 3 shows a nuclear and off-nuclear $`H+K`$ spectrum. Some properties deserve mentioning: * The nuclear spectrum is clearly reddened compared to the off-nuclear one. Reddening towards the nucleus due to the presence of dust is typically found in AGN (Glass & Moorwood, 1985). * The nuclear Pa$`\alpha `$ line shows a broad component with a width of about 4000 km/s arising from the broad line region. This component is not visible in the off-nuclear spectrum. * Around $`1.8\mu `$m the effect of degrading atmospheric transmission is visible. * Stellar absorption features like the NaI $`\lambda 2.206,2.208`$ doublet, CaI $`\lambda 2.263`$, CO(6-3), and CO(2-0) are visible, helping to estimate the stellar content. ### 2.2 Tracing the narrow-line emitting gas The great advantage of integral field spectroscopy is the simultaneous availability of spatial and spectral information. This allows to generate spatial maps of spectral features of interest. The first three panels in Fig. 4 show the recombination lines Pa$`\alpha `$, Br$`\gamma `$, and HeI as tracers of star formation. This emission is extended and, besides the nucleus, is concentrated in a ring like structure at a projected distance of about 500 pc. The morphology roughly follows the continuum contours of Fig. 2. Thus, we could be dealing with a starburst ring embedded in a nuclear bar. There is a plentiful number of examples of starburst rings in galaxies with and without nuclear activity (cf. Smith et al., 1999, and references therein) and which have comparable sizes. Other species like the forbidden transitions \[SiVI\] and \[FeII\] are concentrated on the nucleus in Fig. 4. High energetic nuclear emission is primarily able to penetrate deeply into the interstellar medium and produce regions which are partly ionized and where such transitions can be excited. The \[FeII\] emission is slightly extended, indicating also a connection with supernovae excited emission in the starburst ring. Rotational/vibrational transitions of H<sub>2</sub> can be used to study the dominant excitation mechanisms in the circum-nuclear environment (e.g. Mouri, 1994). The emission lines originate in surfaces of molecular clouds exposed to stellar or nuclear radiation. Three major processes are considered: (1) excitation by X-ray or cosmic rays coming from the AGN; (2) shock excitation either from supernova winds or streaming motion; or (3) UV fluorescence from young stars. The first two processes are of thermal nature, whereas the third is non thermal. 1-0S(1) and the other detectable emission lines are also concentrated on the nucleus. Line ratios are used to estimate the level populations of H<sub>2</sub> in Mrk 609, indicating a thermal origin of the emission with a temperature of about 1900 K (Fig. 5, Zuther et al. in prep.). At the current stage of analysis, excitation due to X-rays also seems to be unimportant, because a number of other H<sub>2</sub> transitions which would be expected in the case of X-ray excitation (Tiné et al., 1997), are not detected. This was already suggested by previous studies of a number of infrared galaxies (e.g. Koorneef & Israel, 1996). Line ratios of hydrogen recombination lines (e.g. Pa$`\alpha `$ and Br$`\gamma `$) can be used to estimate the amount extinction. Assuming case-B recombination we find no significant extinction in the ring like structure. The nuclear ratios also give no significant extinction. At the nucleus, however, this value can be influenced by not-perfect separation of the narrow and broad component of the Pa$`\alpha `$ line. One can further use parts of the spectrum without prominent emission/absorption lines to generate a reddening map. Fig. 6 shows the circum-nuclear reddening using parts of the continuum in the $`H`$\- and the $`K`$-band. The largest values of reddening are found close to/at the nucleus and at the local Pa$`\alpha `$ peaks at the tips of the bar within the putative starburst ring. This indicates the presence of dust and molecular material feeding the AGN or the star formation. ## 3 Summary and conclusions In this contribution we have presented SINFONI science verification observations of the starburst/Seyfert composite galaxy Mrk 609, which has been drawn from a sample of AO-suitable X-ray luminous AGN. AO-assisted integral field spectroscopy in the NIR providing simultaneous spatial and spectroscopic information enables the detailed study of the connection/feedback between star formation and nuclear activity by means of emission/absorption line diagnostics, 2-dimensional morphological, and kinematical analysis. The presented observations indicate the presence of a nuclear bar and associated star formation in a starburst ring. The dominating central excitation mechanism for molecular hydrogen appears to be of thermal origin. Further work on the spectral/spatial analysis will quantify the above statements and incorporation of complementary multi wavelength data (SDSS, ROSAT, CO(1-0); Zuther et al. in prep.) will give a better understanding of the feedback between nuclear activity and host galaxy environment in Mrk 609. Expanding this study to other members of the sample will provide further insights into the nature of the class of starbust/Seyfert composites itself. Acknowledgments This work was supported in part by the Deutsche Forschungsgemeinschaft (DFG) via grant SFB 494.
warning/0506/hep-ph0506067.html	ar5iv	text	# The Decay 𝚺⁺→𝒑⁢ℓ⁺⁢ℓ⁻ within the Standard Model ## I Introduction Three events for the decay mode $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$ have been recently observed by the HyperCP (E871) collaboration Park:2005ek with results that suggest new physics may be needed to explain them. In this paper we re-examine this mode Bergstrom:1987wr within the standard model. There are short- and long-distance contributions to this decay. In the standard model (SM), the leading short-distance contribution comes from the $`Z`$-penguin and box diagrams, as well as the electromagnetic penguin with the photon connected to the dimuon pair Buchalla:1995vs . We find that this contribution yields a branching ratio of order $`10^{12}`$, which is much smaller than the central experimental value of $`8.6\times 10^8`$ reported by HyperCP Park:2005ek . It is well known that the long-distance contribution to the weak radiative mode $`\mathrm{\Sigma }^+p\gamma `$ is much larger than the short-distance contribution. It is therefore also possible to have enhanced long-distance contributions to $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$ via an intermediate virtual photon from $`\mathrm{\Sigma }^+p\gamma `$. We find that the resulting branching ratio is in agreement with the measured value. There is, of course, still the possibility He:1999ik that new physics is responsible for the observed branching ratio of $`\mathrm{\Sigma }^+p\gamma `$ and hence that of $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$. This implies that it is essential to have an up-to-date estimate of the standard-model contributions, on which we concentrate in this work. In Sec. II we update the estimate of the short-distance amplitude. We use the standard effective Hamiltonian for the $`sd\mathrm{}^+\mathrm{}^{}`$ transition Buchalla:1995vs supplemented with hadronic matrix elements for the relevant currents. In Sec. III we study the long-distance contributions mediated by a real or a virtual photon. These can be parameterized by four (complex) gauge-invariant form-factors Bergstrom:1987wr . We determine the imaginary parts of these form factors from unitarity. The real parts of two of the form factors can be reasonably assumed to be constant as a first approximation and can then be extracted from the measured rate and asymmetry parameter for $`\mathrm{\Sigma }^+p\gamma `$ up to a fourfold ambiguity. The real parts of the two remaining form-factors cannot be extracted from experiment at present, and so we estimate them using vector-meson-dominance models. Finally, in Sec. IV we combine all these results to present the predictions for the rates and spectra of the two modes $`\mathrm{\Sigma }^+p\mu ^+\mu ^{},pe^+e^{}`$. Before concluding, we discuss the implications of our analysis for the possibility that new physics could be present in the recent measurement by HyperCP. ## II Short-distance contributions The short-distance effective Hamiltonian responsible for $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ contains contributions originating from the $`Z`$-penguin, box, and electromagnetic-penguin diagrams. It is given by Shifman:1976de ; Buchalla:1995vs $`_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{ud}^{}V_{us}\left[\left(z_{7V}+\tau y_{7V}\right)O_{7V}+\tau y_{7A}O_{7A}\right]+{\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \underset{j}{}}V_{jd}^{}V_{js}c_{7\gamma }^jO_{7\gamma },`$ (1) where $`V_{kl}`$ are the elements of the Cabibbo-Kobayashi-Maskawa (CKM) matrix ckm , $`z`$, $`y`$, and $`c`$ are the Wilson coefficients, $`\tau =V_{td}^{}V_{ts}/\left(V_{ud}^{}V_{us}\right)`$, and $`O_{7V}`$ $`=`$ $`\overline{d}\gamma ^\mu (1\gamma _5)s\overline{\mathrm{}}^{}\gamma _\mu \mathrm{}^+,O_{7A}=\overline{d}\gamma ^\mu (1\gamma _5)s\overline{\mathrm{}}^{}\gamma _\mu \gamma _5\mathrm{}^+,`$ $`O_{7\gamma }`$ $`=`$ $`{\displaystyle \frac{e}{16\pi ^2}}\overline{d}\sigma ^{\mu \nu }F_{\mu \nu }\left[m_s(1+\gamma _5)+m_d(1\gamma _5)\right]s,`$ (2) with $`F_{\mu \nu }`$ being the photon field-strength tensor. The contribution of $`O_{7\gamma }`$ to $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ occurs via the photon converting to a lepton pair. The total short-distance contribution to the $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ amplitude is then given by $`(\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{})=p\mathrm{}^+\mathrm{}^{}\|_{\mathrm{eff}}\|\mathrm{\Sigma }^+`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\{V_{ud}^{}V_{us}[(z_{7V}+\tau y_{7V})p\|\overline{d}\gamma ^\mu (1\gamma _5)s\|\mathrm{\Sigma }^+\overline{\mathrm{}}^{}\gamma _\mu \mathrm{}^++\tau y_{7A}p\|\overline{d}\gamma ^\mu (1\gamma _5)s\|\mathrm{\Sigma }^+\overline{\mathrm{}}^{}\gamma _\mu \gamma _5\mathrm{}^+]`$ $`{\displaystyle \underset{j}{}}V_{jd}^{}V_{js}{\displaystyle \frac{i\alpha c_{7\gamma }^j}{2\pi q^2}}[(m_s+m_d)p\|\overline{d}\sigma ^{\mu \nu }q_\nu s\|\mathrm{\Sigma }^++(m_sm_d)p\|\overline{d}\sigma ^{\mu \nu }q_\nu \gamma _5s\|\mathrm{\Sigma }^+]\overline{\mathrm{}}^{}\gamma _\mu \mathrm{}^+\},`$ where $`q=p_\mathrm{\Sigma }p_p`$. To obtain the corresponding branching ratio, one needs to know the hadronic matrix elements. Employing the leading-order strong Lagrangian in chiral perturbation theory ($`\chi `$PT), given in Eq. (31), we find $`p\|\overline{d}\gamma ^\mu s\|\mathrm{\Sigma }^+=\overline{p}\gamma ^\mu \mathrm{\Sigma },`$ $`p\|\overline{d}\gamma ^\mu \gamma _5s\|\mathrm{\Sigma }^+=(DF)\overline{p}\gamma ^\mu \gamma _5\mathrm{\Sigma },`$ (4) where $`D=0.80`$ and $`F=0.46`$ from fitting to hyperon semileptonic decays, and using quark-model results Donoghue:1992dd we obtain $`p\|\overline{d}\sigma ^{\mu \nu }s\|\mathrm{\Sigma }^+=c_\sigma \overline{p}\sigma ^{\mu \nu }\mathrm{\Sigma },`$ $`p\|\overline{d}\sigma ^{\mu \nu }\gamma _5s\|\mathrm{\Sigma }^+=c_\sigma \overline{p}\sigma ^{\mu \nu }\gamma _5\mathrm{\Sigma },`$ (5) where $`c_\sigma =1/3`$. Furthermore, we adopt the CKM-matrix elements given in Ref. pdg , the typical Wilson coefficients obtained in the literature Shifman:1976de ; Buchalla:1995vs , namely $`z_{7V}=0.046\alpha `$, $`y_{7V}=0.735\alpha `$, $`y_{7A}=0.700\alpha `$ Buchalla:1995vs , and $`c_{7\gamma }^j`$ being dominated by $`c_{7\gamma }^c=0.13`$ Shifman:1976de , and the quark masses $`m_d=9\mathrm{MeV}`$ and $`m_s=120\mathrm{MeV}`$. The resulting branching ratio for $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$ is about $`10^{12}`$, which is way below the observed value. There are uncertainties in the hadronic matrix elements, the Wilson coefficients, and the CKM-matrix elements, but these uncertainties will not change this result by orders of magnitude. We therefore conclude that in the SM the short-distance contribution is too small to explain the HyperCP data on $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$. Now, a large branching ratio for $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ may be related to the large observed branching ratio for $`\mathrm{\Sigma }^+p\gamma `$, compared with their respective short-distance contributions. With only the short-distance contribution to $`\mathrm{\Sigma }^+p\gamma `$ within the SM, the branching ratio is predicted to be much smaller than the experimental value He:1999ik . However, beyond the SM it is possible to have an enhanced short-distance contribution to $`\mathrm{\Sigma }^+p\gamma `$ He:1999ik which would enhance the amplitude for $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$. The origin of the enhancement may be from new interactions such as $`W_L`$-$`W_R`$ mixing in left-right symmetric models and left-right squark mixing in supersymmetric models He:1999ik . These types of interactions have small effects on other related flavor-changing processes such as $`K^0`$-$`\overline{K}^0`$ mixing, but can have large effects on $`\mathrm{\Sigma }^+p\gamma `$ and therefore also on $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$. Thus the observed branching ratio for $`\mathrm{\Sigma }^+p\gamma `$ can be reproduced even if one assumes that there is only the short-distance contribution. More likely, however, the enhancement is due to long-distance contributions within the SM. In the next section we present the most complete estimate possible at present for these long-distance contributions. ## III Long-distance contributions In this section we deal with the contributions to $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ that are mediated by a photon. For a real intermediate photon there are two form factors that can be extracted from the weak radiative hyperon decay $`B_iB_f\gamma `$ and are usually parameterized by the effective Lagrangian $$=\frac{eG_F}{2}\overline{B}_f\left(a+b\gamma _5\right)\sigma ^{\mu \nu }B_iF_{\mu \nu }.$$ (6) The two form factors, $`a`$ and $`b`$, are related to the width and decay distribution of the radiative decay by $`\mathrm{\Gamma }(B_iB_f\gamma )`$ $`=`$ $`{\displaystyle \frac{G_F^2e^2}{\pi }}\left(\|a\|^2+\|b\|^2\right)\omega ^3,`$ (7) $`{\displaystyle \frac{d\mathrm{\Gamma }}{d\mathrm{cos}\theta }}\mathrm{\hspace{0.17em}\hspace{0.17em}1}+\alpha \mathrm{cos}\theta ,\alpha ={\displaystyle \frac{2Re(ab^{})}{\|a\|^2+\|b\|^2}},`$ (8) where $`\omega `$ is the photon energy, and $`\theta `$ is the angle between the spin of $`B_i`$ and the three-momentum of $`B_f`$. The measured values for $`\mathrm{\Sigma }^+p\gamma `$ are pdg $`\mathrm{\Gamma }(\mathrm{\Sigma }^+p\gamma )=(10.1\pm 0.4)\times 10^{15}\mathrm{MeV},\alpha =0.76\pm 0.08.`$ (9) When the photon is a virtual one, there are two additional form-factors, and the total amplitude can be parameterized as $`(B_iB_f\gamma ^{})`$ $`=`$ $`eG_F\overline{B}_f\left[i\sigma ^{\mu \nu }q_\mu (a+b\gamma _5)+(q^2\gamma ^\nu q^\nu \overline{)}q)(c+d\gamma _5)\right]B_i\epsilon _\nu ^{},`$ (10) where $`q`$ is the photon four-momentum. We note that the $`a`$ and $`c`$ ($`b`$ and $`d`$) terms are parity conserving (violating). The corresponding amplitude for $`B_iB_f\mathrm{}^+\mathrm{}^{}`$ is then $`(B_iB_f\mathrm{}^+\mathrm{}^{})`$ $`=`$ $`{\displaystyle \frac{ie^2G_F}{q^2}}\overline{B}_f\left(a+b\gamma _5\right)\sigma _{\mu \nu }q^\mu B_i\overline{\mathrm{}}^{}\gamma ^\nu \mathrm{}^+`$ (11) $`e^2G_F\overline{B}_f\gamma _\mu (c+d\gamma _5)B_i\overline{\mathrm{}}^{}\gamma ^\mu \mathrm{}^+,`$ where now $`q=p_\mathrm{}^++p_{\mathrm{}^{}}.`$ In general $`a`$, $`b`$, $`c`$, and $`d`$ depend on $`q^2`$, and for $`\mathrm{\Sigma }^+p\gamma ^{}`$ the first two are constrained at $`q^2=0`$ by the data in Eq. (9) as $`\|a(0)\|^2+\|b(0)\|^2`$ $`=`$ $`(15.0\pm 0.3)^2\mathrm{MeV}^2,`$ $`\mathrm{Re}\left(a(0)b^{}(0)\right)`$ $`=`$ $`(85.3\pm 9.6)\mathrm{MeV}^2.`$ (12) These form factors are related to the ones in Ref. Bergstrom:1987wr by $`a=\mathrm{\hspace{0.17em}\hspace{0.17em}2}ib_1,b=\mathrm{\hspace{0.17em}\hspace{0.17em}2}ib_2,c={\displaystyle \frac{ia_1}{q^2}},d={\displaystyle \frac{ia_2}{q^2}}.`$ (13) As we will estimate later on, these form factors have fairly mild $`q^2`$-dependence. If they are taken to be constant, by integrating numerically over phase space we can determine the branching ratios of $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ to be, with $`a`$ and $`b`$ in MeV, $`(\mathrm{\Sigma }^+p\mu ^+\mu ^{})`$ $`=`$ $`\left[2.00\left(\|a\|^2+\|b\|^2\right)1.60\left(\|a\|^2\|b\|^2\right)\right]\times 10^{10}`$ (14a) $`+\left(1.05\|c\|^2+18.2\|d\|^2\right)\times 10^6`$ $`+\left[0.29\mathrm{Re}(ac^{})16.1\mathrm{Re}(bd^{})\right]\times 10^8,`$ $`(\mathrm{\Sigma }^+pe^+e^{})`$ $`=`$ $`\left[4.22\left(\|a\|^2+\|b\|^2\right)0.21\left(\|a\|^2\|b\|^2\right)\right]\times 10^8`$ (14b) $`+\left(5.38\|c\|^2+15.9\|d\|^2\right)\times 10^5`$ $`+\left[1.51\mathrm{Re}(ac^{})21.1\mathrm{Re}(bd^{})\right]\times 10^7.`$ If the form factors have $`q^2`$-dependence, the expression is different, and the rate should be calculated with the formula which we give in Appendix A. ### III.1 Imaginary parts of the form factors from unitarity The form factors which contribute to the weak radiative hyperon decays have been studied in chiral perturbation theory Neufeld:1992hb ; Jenkins:1992ab ; Bos:1996ig . The imaginary parts of $`a`$ and $`b`$ for $`\mathrm{\Sigma }^+p\gamma `$ have been determined from unitarity with different results in the literature. Neufeld Neufeld:1992hb employed relativistic baryon $`\chi `$PT to find, for $`q^2=0`$, $`\mathrm{Im}a(0)=\mathrm{\hspace{0.17em}\hspace{0.17em}2.60}\mathrm{MeV},\mathrm{Im}b(0)=1.46\mathrm{MeV}`$ (15) in the notation of Eq. (6), whereas Jenkins et al. Jenkins:1992ab using the heavy-baryon formulation obtained $`\mathrm{Im}a(0)=\mathrm{\hspace{0.17em}\hspace{0.17em}6.18}\mathrm{MeV},\mathrm{Im}b(0)=0.53\mathrm{MeV}.`$ (16) Because of this disagreement, and since we also need the imaginary parts of the form factors $`c`$ and $`d`$, we repeat here the unitarity calculation employing both the relativistic and heavy baryon approaches. Our strategy to derive the imaginary parts of the four form-factors in Eq. (11) from unitarity is illustrated in Fig. 1. As the figure shows, these imaginary parts can be determined from the amplitudes for the weak nonleptonic decays $`\mathrm{\Sigma }^+p\pi ^0`$ and $`\mathrm{\Sigma }^+n\pi ^+`$ (the vertex indicated by a square in Fig. 1) as well as the reactions $`N\pi N\gamma ^{}`$ (the vertex indicated by a blob in Fig. 1). The weak decays have been measured pdg , and we express their amplitudes as<sup>1</sup><sup>1</sup>1 We have taken the nonzero elements of $`\gamma _5`$ to be positive. $`(\mathrm{\Sigma }^+N\pi )`$ $`=`$ $`iG_Fm_{\pi ^+}^2\overline{N}\left(A_{N\pi }B_{N\pi }\gamma _5\right)\mathrm{\Sigma },`$ (17) where $`A_{n\pi ^+}=\mathrm{\hspace{0.17em}\hspace{0.17em}0.06},`$ $`B_{n\pi ^+}=\mathrm{\hspace{0.17em}\hspace{0.17em}18.53},`$ $`A_{p\pi ^0}=1.43,`$ $`B_{p\pi ^0}=\mathrm{\hspace{0.17em}\hspace{0.17em}11.74}.`$ (18) Following Refs. Neufeld:1992hb ; Jenkins:1992ab , we adopt the $`N\pi p\gamma ^{}`$ amplitudes derived in lowest-order $`\chi `$PT. We present the details of our unitarity calculation in Appendix B. The results in the relativistic and heavy baryon approaches are given in Eqs. (32) and (38), respectively. In Fig. 2 we display the two sets of form factors for $`\mathrm{\hspace{0.17em}0}q^2(m_\mathrm{\Sigma }m_N)^2`$. We note that, although only the $`\mathrm{\Sigma }^+n\pi ^+`$ transition contributes to the heavy-baryon form-factors at leading order, the sizable difference between the $`\mathrm{Im}a`$, or $`\mathrm{Im}c`$, curves arises mainly from relativistic corrections, which reduce the heavy-baryon numbers by about 50%. On the other hand, the difference between the $`\mathrm{Im}b`$, or $`\mathrm{Im}d`$, curves is due not only to relativistic corrections, but also to $`A_{n\pi ^+}`$ being much smaller than $`A_{p\pi ^0}`$. To compare with the numbers in Eqs. (15) and (16) calculated in earlier work, we find from the relativistic formulas in Eq. (32) $`\mathrm{Im}a(0)=\mathrm{\hspace{0.17em}\hspace{0.17em}2.84}\mathrm{MeV},\mathrm{Im}b(0)=1.83\mathrm{MeV},`$ (19) and from the heavy-baryon results in Eq. (38) $`\mathrm{Im}a(0)=\mathrm{\hspace{0.17em}\hspace{0.17em}6.84}\mathrm{MeV},\mathrm{Im}b(0)=0.54\mathrm{MeV}.`$ (20) Thus our relativistic results are close to those in Eq. (15), from Ref. Neufeld:1992hb , and our heavy-baryon numbers to those in Eq. (16), from Ref. Jenkins:1992ab .<sup>2</sup><sup>2</sup>2 Our heavy-baryon expressions for $`\mathrm{Im}a(0)`$ and $`\mathrm{Im}b(0)`$ are identical to those in Ref. Jenkins:1992ab , except that their $`\mathrm{Im}a(0)`$ formula has one of the overall factors of $`1/(m_\mathrm{\Sigma }m_N)`$ apparently coming from their approximating $`[(m_\mathrm{\Sigma }m_N)^2m_\pi ^2]^{1/2}`$ as $`m_\mathrm{\Sigma }m_N`$. This is the main reason for the value of $`\mathrm{Im}a(0)`$ in Eq. (16) being smaller than that in Eq. (20). These two sets of numbers are different for the reasons mentioned in the preceding paragraph. ### III.2 Real parts of the form factors The real parts of the form factors cannot be completely predicted at present from experimental input alone. For $`\mathrm{Re}a(q^2)`$ and $`\mathrm{Re}b(q^2)`$, the values at $`q^2=0`$ can be extracted from Eq. (III) after using Eq. (19) or (20) for the imaginary parts. Thus the relativistic numbers in Eq. (19) lead to the four sets of solutions $`\mathrm{Re}a(0)=\pm 13.3\mathrm{MeV},`$ $`\mathrm{Re}b(0)=6.0\mathrm{MeV},`$ $`\mathrm{Re}a(0)=\pm 6.0\mathrm{MeV},`$ $`\mathrm{Re}b(0)=13.3\mathrm{MeV},`$ (21) while the heavy-baryon results in Eq. (20) imply $`\mathrm{Re}a(0)=\pm 11.1\mathrm{MeV},`$ $`\mathrm{Re}b(0)=7.3\mathrm{MeV},`$ $`\mathrm{Re}a(0)=\pm 7.3\mathrm{MeV},`$ $`\mathrm{Re}b(0)=11.1\mathrm{MeV}.`$ (22) Since these numbers still cannot be predicted reliably within the framework of $`\chi `$PT Neufeld:1992hb ; Jenkins:1992ab , we will assume that $`\mathrm{Re}a(q^2)=\mathrm{Re}a(0),`$ $`\mathrm{Re}b(q^2)=\mathrm{Re}b(0),`$ (23) where the $`q^2=0`$ values are those in Eqs. (III.2) and (III.2) in the respective approaches. This assumption is also reasonable in view of the fairly mild $`q^2`$-dependence of the imaginary parts seen in Fig. 2, and of the real parts of $`c`$ and $`d`$ below. In predicting the $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ rates in the following section, we will use the 8 sets of possible solutions in Eqs. (III.2) and (III.2). The real parts of $`c`$ and $`d`$ cannot be extracted from experiment at present. Our interest here, however, is in predicting the SM contribution, and therefore we need to estimate them. To do so, we employ a vector-meson-dominance assumption, presenting the details in Appendix C. The results for $`\mathrm{Re}c(q^2)`$ and $`\mathrm{Re}d(q^2)`$ are given in Eqs. (41) and (43), respectively. In Fig. 3 we display the two form factors for $`\mathrm{\hspace{0.17em}0}q^2(m_\mathrm{\Sigma }m_N)^2`$. We can see from Figs. 2 and 3 that $`c`$ is dominated by its imaginary part, but that $`d`$ is mostly real. ## IV Results and conclusions We can now evaluate the rates and spectra of $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ resulting from the various standard-model contributions. Since the short-distance contributions discussed in Sec. II are very small, we shall neglect them. Consequently, the rates are determined by the various form factors in $`\mathrm{\Sigma }^+p\gamma ^{}`$ calculated in the preceding section and applied in Eq. (27). In Table 1, we have collected the branching ratios of $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$ and $`\mathrm{\Sigma }^+pe^+e^{}`$ corresponding to the 8 sets of solutions in Eqs. (III.2) and (III.2), under the assumption of Eq. (23) for $`\mathrm{Re}a`$ and $`\mathrm{Re}b`$. The real parts of $`c`$ and $`d`$ in Eqs. (41) and (43) are used in all the unbracketed branching ratios. For the imaginary parts of the form factors, the expressions in Eq. (32) \[Eq. (38)\] contribute to the unbracketed branching ratios in the upper (lower) half of this table. Within each pair of square brackets, the first number is the branching ratio obtained without contributions from both $`c`$ and $`d`$, whereas the second number is the branching ratio calculated with only the real parts of all the form factors. In Fig. 4 we show the invariant-mass distributions of the $`\mu ^+\mu ^{}`$ pair, with $`M_{\mu \mu }=\sqrt{q^2}`$, that correspond to the smallest and largest rates of $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$ listed in Table 1 for both the relativistic baryon \[(a) and (b)\] and heavy baryon \[(c) and (d)\] cases. For $`\mathrm{\Sigma }^+pe^+e^{}`$, the mass distributions of the $`e^+e^{}`$ pair, two of which are displayed in Fig. 5, differ very little from each other and are strongly peaked at low $`M_{ee}=\sqrt{q^2}`$. Also shown in the figures are the distributions obtained with $`c=d=0`$ (dashed curves), as well as those without contributions from the imaginary parts of all the form factors (dotted curves). We can see from Table 1, Fig. 4, and Fig. 5 that the effect of the $`c`$ and $`d`$ contributions on the total rates can be up to nearly 40% in $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$, but it is much smaller in $`\mathrm{\Sigma }^+pe^+e^{}`$. Furthermore, the contributions of the imaginary parts of the form factors can be as large as 35% to the $`p\mu ^+\mu ^{}`$ rate and roughly 20% to the $`pe^+e^{}`$ rate. This implies that a careful analysis of experimental results, especially in the case of $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$, should take into account the imaginary parts of the form factors. For $`\mathrm{\Sigma }^+p\mu ^+\mu ^{}`$, HyperCP measured the branching ratio to be $`\left(8.6_{5.4}^{+6.6}\pm 5.5\right)\times 10^8`$ Park:2005ek . It is evident that all the predictions in Table 1 for the $`p\mu ^+\mu ^{}`$ mode corresponding to the different sets of form factors fall within the experimental range. For $`\mathrm{\Sigma }^+pe^+e^{}`$, the branching ratio can be inferred from the experimental results given in Ref. ang , which reported the width ratio $`\mathrm{\Gamma }(\mathrm{\Sigma }^+pe^+e^{})/\mathrm{\Gamma }(\mathrm{\Sigma }^+p\pi ^0)=(1.5\pm 0.9)\times 10^5`$ and interpreted the observed events as proceeding from $`\mathrm{\Sigma }^+p\gamma ^{}`$, based on the very low invariant-masses of the $`e^+e^{}`$ pair.<sup>3</sup><sup>3</sup>3We note that the upper limit of $`7\times 10^6`$ quoted in Ref. pdg and obtained in Ref. ang is for the presence of weak neutral currents in $`\mathrm{\Sigma }^+pe^+e^{}`$ and not for the branching ratio of this mode. This number, in conjunction with the current data on $`\mathrm{\Sigma }^+p\pi ^0`$ pdg , translates into $`(\mathrm{\Sigma }^+pe^+e^{})=(7.7\pm 4.6)\times 10^6`$. Clearly, the results for the $`pe^+e^{}`$ mode in Table 1 are well within the experimentally allowed range. Based on the numbers in Table 1, we may then conclude that within the standard model $`\begin{array}{c}1.6\times 10^8\left(\mathrm{\Sigma }^+p\mu ^+\mu ^{}\right)\mathrm{\hspace{0.17em}\hspace{0.17em}9.0}\times 10^8,\\ 9.1\times 10^6\left(\mathrm{\Sigma }^+pe^+e^{}\right)\mathrm{\hspace{0.17em}\hspace{0.17em}10.1}\times 10^6.\end{array}`$ (26) The agreement above between the predicted and observed rates of $`\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{}`$ indicates that these decays are dominated by long-distance contributions. However, the predicted range for $`(\mathrm{\Sigma }^+p\mu ^+\mu ^{})`$ is sufficiently wide that we cannot rule out the possibility of a new-physics contribution of the type suggested by HyperCP Park:2005ek . Motivated by the narrow distribution of dimuon masses of the events they observed, they proposed that the decay could proceed via a new intermediate particle of mass $``$ 214 MeV, with a branching ratio of $`\left(3.1_{1.9}^{+2.4}\pm 1.5\right)\times 10^8`$ Park:2005ek . For this hypothesis to be realized, however, the new physics would have to dominate the decay. It will be interesting to see if this hypothesis will be confirmed by future measurements. Finally, we observe that the smaller numbers $`(\mathrm{\Sigma }^+p\mu ^+\mu ^{})2\times 10^8`$ in Table 1 correspond to the mass distributions peaking at lower masses, $`M_{\mu \mu }220\mathrm{MeV}`$, in Fig. 4. It is perhaps not coincidental that these numbers are similar to the branching ratio and new-particle mass, respectively, in the HyperCP hypothesis above. This may be another indication that it is not necessary to invoke new physics to explain the HyperCP results. ###### Acknowledgements. We thank HyangKyu Park for conversations. The work of X.G.H. was supported in part by the National Science Council under NSC grants. The work of G.V. was supported in part by DOE under contract number DE-FG02-01ER41155. ## Appendix A Differential rate of $`𝚺^\mathbf{+}\mathbf{}𝒑\mathbf{}^\mathbf{+}\mathbf{}^{\mathbf{}}`$ If the form factors have $`q^2`$-dependence, before integrating over phase space to obtain the branching ratio we should use $`{\displaystyle \frac{d\mathrm{\Gamma }(\mathrm{\Sigma }^+p\mathrm{}^+\mathrm{}^{})}{dq^2dt}}={\displaystyle \frac{\alpha ^2G_F^2}{4\pi m_\mathrm{\Sigma }^3}}`$ (27) $`\times `$ $`\{[(2m_l^2+q^2)((m_pm_\mathrm{\Sigma })^2q^2)(m_\mathrm{\Sigma }+m_p)^2+2q^2f(m_p,m_\mathrm{\Sigma },m_l,q^2,t)]{\displaystyle \frac{\|a\|^2}{q^4}}`$ $`+\left[(2m_l^2+q^2)((m_p+m_\mathrm{\Sigma })^2q^2)(m_\mathrm{\Sigma }m_p)^2+2q^2f(m_p,m_\mathrm{\Sigma },m_l,q^2,t)\right]{\displaystyle \frac{\|b\|^2}{q^4}}`$ $`+\left[(2m_l^2+q^2)((m_pm_\mathrm{\Sigma })^2q^2)2f(m_p,m_\mathrm{\Sigma },m_l,q^2,t)\right]\|c\|^2`$ $`+\left[(2m_l^2+q^2)((m_p+m_\mathrm{\Sigma })^2q^2)2f(m_p,m_\mathrm{\Sigma },m_l,q^2,t)\right]\|d\|^2`$ $`+\mathrm{\hspace{0.17em}\hspace{0.17em}2}(m_\mathrm{\Sigma }+m_p)(2m_l^2+q^2)[(m_pm_\mathrm{\Sigma })^2q^2)]{\displaystyle \frac{\mathrm{Re}(ac^{})}{q^2}}`$ $`2(m_\mathrm{\Sigma }m_p)(2m_l^2+q^2)[(m_p+m_\mathrm{\Sigma })^2q^2]{\displaystyle \frac{\mathrm{Re}(bd^{})}{q^2}}\},`$ where $`t=(p_\mathrm{\Sigma }p_{\mathrm{}^{}})^2`$ and $`f(m_p,m_\mathrm{\Sigma },m_l,q^2,t)=m_l^4+(m_p^2+m_\mathrm{\Sigma }^2q^22t)m_l^2+m_p^2m_\mathrm{\Sigma }^2(m_p^2+m_\mathrm{\Sigma }^2)t+(q^2+t)t,`$ with the integration intervals given by $`\begin{array}{c}t_{\mathrm{max},\mathrm{min}}=\frac{1}{2}\left[m_\mathrm{\Sigma }^2+m_p^2+2m_l^2q^2\pm \sqrt{1{\displaystyle \frac{4m_l^2}{q^2}}}\sqrt{(m_\mathrm{\Sigma }^2m_p^2q^2)^24m_p^2q^2}\right],\\ q_{\mathrm{min}}^2=\mathrm{\hspace{0.17em}\hspace{0.17em}4}m_l^2,q_{\mathrm{max}}^2=(m_\mathrm{\Sigma }m_p)^2.\end{array}`$ (30) It is worth mentioning that, since the form factors belong to the $`\mathrm{\Sigma }^+p\gamma ^{}`$ amplitude, they do not depend on $`t`$. ## Appendix B Imaginary parts of form factors in $`𝝌`$PT The chiral Lagrangian for the interactions of the lowest-lying mesons and baryons is written down in terms of the lightest meson-octet and baryon-octet fields, which are collected into $`3\times 3`$ matrices $`\phi `$ and $`B`$, respectively Bijnens:1985kj . The mesons enter through the exponential $`\mathrm{\Sigma }=\xi ^2=\mathrm{exp}(\mathrm{i}\phi /f),`$ where $`f=f_\pi =92.4\mathrm{MeV}`$ is the pion decay constant. In the relativistic baryon $`\chi `$PT, the lowest-order strong Lagrangian is given by Bijnens:1985kj $`_\mathrm{s}`$ $`=`$ $`\overline{B}i\gamma ^\mu \left(_\mu B+[𝒱_\mu ,B]\right)+m_0\overline{B}B+D\overline{B}\gamma ^\mu \gamma _5\{𝒜_\mu ,B\}+F\overline{B}\gamma ^\mu \gamma _5[𝒜_\mu ,B],`$ (31) where $`\mathrm{}\mathrm{Tr}(\mathrm{})`$ in flavor space, $`m_0`$ is the baryon mass in the chiral limit, $`𝒱^\mu =\frac{1}{2}\left(\xi ^\mu \xi ^{}+\xi ^{}^\mu \xi \right)+\frac{i}{2}eA^\mu \left(\xi ^{}Q\xi +\xi Q\xi ^{}\right),`$ and $`𝒜^\mu =\frac{i}{2}\left(\xi ^\mu \xi ^{}\xi ^{}^\mu \xi \right)+\frac{1}{2}eA^\mu \left(\xi ^{}Q\xi \xi Q\xi ^{}\right),`$ with $`A^\mu `$ being the photon field and $`Q=\mathrm{diag}(2,1,1)/3`$ the quark-charge matrix.<sup>4</sup><sup>4</sup>4 Under a chiral transformation, $`\overline{B}U\overline{B}U^{}`$, $`BUBU^{}`$, $`𝒱^\mu U𝒱^\mu U^{}+i^\mu UU^{}`$, and $`𝒜^\mu U𝒜^\mu U^{}`$, where $`U`$ is defined by $`\xi L\xi U^{}=U\xi R^{}`$. The parameters $`D`$ and $`F`$ will enter our results below only through the combination $`D+F=1.26.`$ From $`_\mathrm{s}`$ we derive two sets of diagrams, shown in Fig. 6, which represent the $`N\pi p\gamma ^{}`$ reactions involved in the unitarity calculation of the imaginary parts of the form factors $`a`$, $`b`$, $`c`$, and $`d`$. It then follows from Fig. 1 that the first set of diagrams is associated with the weak transition $`\mathrm{\Sigma }^+n\pi ^+`$, and the second with $`\mathrm{\Sigma }^+p\pi ^0`$. Consequently, we express our results as $`\mathrm{Im}`$ $`=`$ $`{\displaystyle \frac{(D+F)m_{\pi ^+}^2}{8\sqrt{2}\pi f_\pi }}\left(\stackrel{~}{}_++{\displaystyle \frac{\stackrel{~}{}_0}{\sqrt{2}}}\right)\text{for }=a,b,c,d,`$ (32) where $`\stackrel{~}{}_+`$ $`\left(\stackrel{~}{}_0\right)`$ comes from the $`n\pi ^+`$ $`(p\pi ^0)`$ contribution, and write them in terms of the weak amplitudes $`A_+=A_{n\pi ^+}`$, $`A_0=A_{p\pi ^0}`$, $`B_+=B_{n\pi ^+}`$, and $`B_0=B_{p\pi ^0}`$ given in Eq. (III.1). Working in the $`\mathrm{\Sigma }^+`$ rest-frame, which implies that the energies and momenta of the photon and proton in the final state and of the pion in the intermediate are fixed by kinematics, we define $`z_+=\left({\displaystyle \frac{2E_\pi E_\gamma 2\|𝒑_\pi \|\|𝒑_\gamma \|q^2}{2E_\pi E_\gamma +2\|𝒑_\pi \|\|𝒑_\gamma \|q^2}}\right),`$ $`z_0=\left({\displaystyle \frac{2E_\pi E_p2\|𝒑_\pi \|\|𝒑_p\|m_\pi ^2}{2E_\pi E_p+2\|𝒑_\pi \|\|𝒑_p\|m_\pi ^2}}\right).`$ (33) The expression for $`\stackrel{~}{}`$ from each set of diagrams can then be written as $`\stackrel{~}{a}_{+,0}`$ $`=`$ $`{\displaystyle \frac{B_{+,0}m_N}{2m_\mathrm{\Sigma }^2\|𝒑_\gamma \|}}{\displaystyle \frac{\left[2\|𝒑_\pi \|\|𝒑_\gamma \|f_{+,0}^{(a)}+\mathrm{ln}(z_{+,0})g_{+,0}^{(a)}\right]}{\left[(m_\mathrm{\Sigma }m_N)^2q^2\right]\left[(m_\mathrm{\Sigma }+m_N)^2q^2\right]^2}},`$ $`\stackrel{~}{b}_{+,0}`$ $`=`$ $`{\displaystyle \frac{A_{+,0}m_N}{2m_\mathrm{\Sigma }^2\|𝒑_\gamma \|}}{\displaystyle \frac{\left[2\|𝒑_\pi \|\|𝒑_\gamma \|f_{+,0}^{(b)}+\mathrm{ln}(z_{+,0})g_{+,0}^{(b)}\right]}{\left[(m_\mathrm{\Sigma }m_N)^2q^2\right]^2\left[(m_\mathrm{\Sigma }+m_N)^2q^2\right]}},`$ $`\stackrel{~}{c}_{+,0}`$ $`=`$ $`{\displaystyle \frac{B_{+,0}m_N}{2m_\mathrm{\Sigma }^2\|𝒑_\gamma \|(m_Nm_\mathrm{\Sigma })}}{\displaystyle \frac{\left[2\|𝒑_\pi \|\|𝒑_\gamma \|f_{+,0}^{(c)}+\mathrm{ln}(z_{+,0})g_{+,0}^{(c)}\right]}{\left[(m_\mathrm{\Sigma }m_N)^2q^2\right]\left[(m_\mathrm{\Sigma }+m_N)^2q^2\right]^2}},`$ $`\stackrel{~}{d}_{+,0}`$ $`=`$ $`{\displaystyle \frac{A_{+,0}m_N}{2m_\mathrm{\Sigma }^2\|𝒑_\gamma \|(m_\mathrm{\Sigma }+m_N)}}{\displaystyle \frac{\left[2\|𝒑_\pi \|\|𝒑_\gamma \|f_{+,0}^{(d)}+\mathrm{ln}(z_{+,0})g_{+,0}^{(d)}\right]}{\left[(m_\mathrm{\Sigma }m_N)^2q^2\right]^2\left[(m_\mathrm{\Sigma }+m_N)^2q^2\right]}},`$ (34) where $`f_+^{(a)}`$ $`=`$ $`m_Nm_\mathrm{\Sigma }^5+\left(q^2+2m_\pi ^2+m_N^2\right)m_\mathrm{\Sigma }^4m_N\left(3q^23m_\pi ^2+2m_N^2\right)m_\mathrm{\Sigma }^3`$ $`\left(q^45m_\pi ^2q^2+2m_N^4+\left(q^2+m_\pi ^2\right)m_N^2\right)m_\mathrm{\Sigma }^2`$ $`+m_N\left(m_N^2q^2\right)\left(2q^23m_\pi ^2+m_N^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\left(m_N^2m_\pi ^2\right),`$ $`g_+^{(a)}`$ $`=`$ $`m_\mathrm{\Sigma }(m_Nq^6+(m_N(2m_Nm_\mathrm{\Sigma })(m_N+m_\mathrm{\Sigma })m_\pi ^2(3m_N+m_\mathrm{\Sigma }))q^4`$ $`+m_\pi ^2(3m_\pi ^24m_N^2)(m_N+m_\mathrm{\Sigma })q^2+m_\pi ^2(m_Nm_\mathrm{\Sigma })^2(m_N+m_\mathrm{\Sigma })^3),`$ $`f_0^{(a)}`$ $`=`$ $`3m_Nm_\mathrm{\Sigma }^5\left(q^22m_\pi ^23m_N^2\right)m_\mathrm{\Sigma }^4m_N\left(4m_N^23\left(q^2+m_\pi ^2\right)\right)m_\mathrm{\Sigma }^3`$ $`+\left(q^4+5m_\pi ^2q^24m_N^4\left(q^2+m_\pi ^2\right)m_N^2\right)m_\mathrm{\Sigma }^2`$ $`+m_N\left(m_N^2q^2\right)\left(2q^23m_\pi ^2+m_N^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\left(m_N^2m_\pi ^2\right),`$ $`g_0^{(a)}`$ $`=`$ $`m_\mathrm{\Sigma }(2m_\pi ^2m_\mathrm{\Sigma }q^4(m_N+m_\mathrm{\Sigma })(3m_\pi ^42(3m_N^22m_\mathrm{\Sigma }m_N+m_\mathrm{\Sigma }^2)m_\pi ^2`$ (35a) $`+m_N(m_Nm_\mathrm{\Sigma })^2(3m_N+m_\mathrm{\Sigma }))q^2+m_N(m_Nm_\mathrm{\Sigma })^2m_\mathrm{\Sigma }(m_N+m_\mathrm{\Sigma })^3),`$ $`f_+^{(b)}`$ $`=`$ $`m_Nm_\mathrm{\Sigma }^5\left(q^2+2m_\pi ^2+m_N^2\right)m_\mathrm{\Sigma }^4m_N\left(3q^23m_\pi ^2+2m_N^2\right)m_\mathrm{\Sigma }^3`$ $`+\left(q^45m_\pi ^2q^2+2m_N^4+\left(q^2+m_\pi ^2\right)m_N^2\right)m_\mathrm{\Sigma }^2`$ $`+m_N\left(m_N^2q^2\right)\left(2q^23m_\pi ^2+m_N^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\left(m_\pi ^2m_N^2\right),`$ $`g_+^{(b)}`$ $`=`$ $`m_\mathrm{\Sigma }(m_Nq^6+((3m_Nm_\mathrm{\Sigma })m_\pi ^2+m_N(2m_N^2+m_\mathrm{\Sigma }m_N+m_\mathrm{\Sigma }^2))q^4`$ $`m_\pi ^2(3m_\pi ^24m_N^2)(m_Nm_\mathrm{\Sigma })q^2m_\pi ^2(m_Nm_\mathrm{\Sigma })^3(m_N+m_\mathrm{\Sigma })^2),`$ $`f_0^{(b)}`$ $`=`$ $`3m_Nm_\mathrm{\Sigma }^5+\left(q^22m_\pi ^23m_N^2\right)m_\mathrm{\Sigma }^4+m_N\left(3\left(q^2+m_\pi ^2\right)4m_N^2\right)m_\mathrm{\Sigma }^3`$ $`\left(q^4+5m_\pi ^2q^24m_N^4\left(q^2+m_\pi ^2\right)m_N^2\right)m_\mathrm{\Sigma }^2`$ $`+m_N\left(m_N^2q^2\right)\left(2q^23m_\pi ^2+m_N^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\left(m_\pi ^2m_N^2\right),`$ $`g_0^{(b)}`$ $`=`$ $`m_\mathrm{\Sigma }(2m_\pi ^2m_\mathrm{\Sigma }q^4+(m_Nm_\mathrm{\Sigma })(3m_\pi ^42(3m_N^2+2m_\mathrm{\Sigma }m_N+m_\mathrm{\Sigma }^2)m_\pi ^2`$ (35b) $`+m_N(3m_Nm_\mathrm{\Sigma })(m_N+m_\mathrm{\Sigma })^2)q^2+m_N(m_Nm_\mathrm{\Sigma })^3m_\mathrm{\Sigma }(m_N+m_\mathrm{\Sigma })^2),`$ $`f_+^{(c)}`$ $`=`$ $`m_\pi ^2\left(8m_\mathrm{\Sigma }^4+5m_Nm_\mathrm{\Sigma }^3\left(3q^2+m_N^2\right)m_\mathrm{\Sigma }^2+3m_N\left(m_N^2q^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\right)`$ $`(m_Nm_\mathrm{\Sigma })(m_Nm_\mathrm{\Sigma }^4+(q^22m_N^2)m_\mathrm{\Sigma }^34q^2m_Nm_\mathrm{\Sigma }^2`$ $`(q^4+m_N^2q^22m_N^4)m_\mathrm{\Sigma }+m_N(q^2m_N^2)^2),`$ $`g_+^{(c)}`$ $`=`$ $`(m_Nm_\mathrm{\Sigma })m_\mathrm{\Sigma }(m_N(2m_N+m_\mathrm{\Sigma })q^4+(m_\pi ^42(3m_N^2+2m_\mathrm{\Sigma }m_N+m_\mathrm{\Sigma }^2)m_\pi ^2`$ $`+m_N(m_Nm_\mathrm{\Sigma })(m_N+m_\mathrm{\Sigma })^2)q^2+2m_\pi ^2(m_N+m_\mathrm{\Sigma })^2(m_\pi ^2+m_\mathrm{\Sigma }(m_\mathrm{\Sigma }m_N))),`$ $`f_0^{(c)}`$ $`=`$ $`m_\pi ^2\left(8m_\mathrm{\Sigma }^4+5m_Nm_\mathrm{\Sigma }^3\left(3q^2+m_N^2\right)m_\mathrm{\Sigma }^2+3m_N\left(m_N^2q^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\right)`$ $`(m_Nm_\mathrm{\Sigma })^2(2m_\mathrm{\Sigma }^4m_Nm_\mathrm{\Sigma }^3(3q^2+m_N^2)m_\mathrm{\Sigma }^2+3m_N(m_N^2q^2)m_\mathrm{\Sigma }`$ $`+(q^2m_N^2)^2),`$ $`g_0^{(c)}`$ $`=`$ $`(m_Nm_\mathrm{\Sigma })m_\mathrm{\Sigma }((m_\pi ^4(2m_N^24m_\mathrm{\Sigma }m_N2m_\mathrm{\Sigma }^2)m_\pi ^2+m_N(m_Nm_\mathrm{\Sigma })^2(m_N+m_\mathrm{\Sigma }))q^2`$ $`+(m_N+m_\mathrm{\Sigma })^2(2m_\pi ^42(2m_N^2m_\mathrm{\Sigma }m_N+m_\mathrm{\Sigma }^2)m_\pi ^2+m_N(m_Nm_\mathrm{\Sigma })^2(2m_Nm_\mathrm{\Sigma }))),`$ $`f_+^{(d)}`$ $`=`$ $`(m_N+m_\mathrm{\Sigma })(m_Nm_\mathrm{\Sigma }^4(q^22m_N^2)m_\mathrm{\Sigma }^34q^2m_Nm_\mathrm{\Sigma }^2`$ $`+(q^4+m_N^2q^22m_N^4)m_\mathrm{\Sigma }+m_N(q^2m_N^2)^2)`$ $`m_\pi ^2\left(8m_\mathrm{\Sigma }^45m_Nm_\mathrm{\Sigma }^3\left(3q^2+m_N^2\right)m_\mathrm{\Sigma }^2+3m_N\left(q^2m_N^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\right),`$ $`g_+^{(d)}`$ $`=`$ $`m_\mathrm{\Sigma }(m_N+m_\mathrm{\Sigma })(m_N(2m_Nm_\mathrm{\Sigma })q^4(m_\pi ^42(3m_N^22m_\mathrm{\Sigma }m_N+m_\mathrm{\Sigma }^2)m_\pi ^2`$ $`+m_N(m_Nm_\mathrm{\Sigma })^2(m_N+m_\mathrm{\Sigma }))q^22m_\pi ^2(m_Nm_\mathrm{\Sigma })^2(m_\pi ^2+m_\mathrm{\Sigma }(m_N+m_\mathrm{\Sigma }))),`$ $`f_0^{(d)}`$ $`=`$ $`\left(m_N+m_\mathrm{\Sigma }\right)^2\left(2m_\mathrm{\Sigma }^4+m_Nm_\mathrm{\Sigma }^3\left(3q^2+m_N^2\right)m_\mathrm{\Sigma }^2+3m_N\left(q^2m_N^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\right)`$ $`m_\pi ^2\left(8m_\mathrm{\Sigma }^45m_Nm_\mathrm{\Sigma }^3\left(3q^2+m_N^2\right)m_\mathrm{\Sigma }^2+3m_N\left(q^2m_N^2\right)m_\mathrm{\Sigma }+\left(q^2m_N^2\right)^2\right),`$ $`g_0^{(d)}`$ $`=`$ $`m_\mathrm{\Sigma }(m_N+m_\mathrm{\Sigma })((m_\pi ^42(m_N^2+2m_\mathrm{\Sigma }m_Nm_\mathrm{\Sigma }^2)m_\pi ^2+m_N(m_Nm_\mathrm{\Sigma })(m_N+m_\mathrm{\Sigma })^2)q^2`$ (35d) $`+(m_Nm_\mathrm{\Sigma })^2(2m_\pi ^42(2m_N^2+m_\mathrm{\Sigma }m_N+m_\mathrm{\Sigma }^2)m_\pi ^2`$ $`+m_N(m_N+m_\mathrm{\Sigma })^2(2m_N+m_\mathrm{\Sigma }))).`$ In our numerical computations, $`m_\mathrm{\Sigma }=m_{\mathrm{\Sigma }^+},`$ $`m_N=\frac{1}{2}\left(m_p+m_n\right),`$ $`m_\pi =\frac{1}{3}\left(2m_{\pi ^+}+m_{\pi ^0}\right),`$ the numbers being from Ref. pdg . In heavy baryon $`\chi `$PT Jenkins:1991ne , the relevant Lagrangian can be found in Ref. Jenkins:1992ab , and the weak radiative and nonleptonic amplitudes in Eqs. (10) and (17) become, respectively, $`(B_iB_f\gamma ^{})`$ $`=`$ $`eG_\mathrm{F}\overline{B}_f\left[2\left(SqS^\mu S^\mu Sq\right)a+2\left(Sqv^\mu S^\mu vq\right)b\right]B_i\epsilon _\mu ^{}`$ (36) $`eG_\mathrm{F}\overline{B}_f\left[\left(q^2v^\mu q^\mu vq\right)c+2\left(q^2S^\mu q^\mu Sq\right)d\right]B_i\epsilon _\mu ^{},`$ $`(\mathrm{\Sigma }^+N\pi )`$ $`=`$ $`iG_Fm_{\pi ^+}^2\overline{N}\left(A_{N\pi }+2Sp_\pi {\displaystyle \frac{B_{N\pi }}{2m_\mathrm{\Sigma }}}\right)\mathrm{\Sigma },`$ (37) where $`v`$ is the baryon four-velocity and $`S`$ is the baryon spin operator. Following Ref. Jenkins:1992ab , to obtain the imaginary parts of the form factors we evaluate the loop diagrams displayed in Fig. 7. In the heavy-baryon approach, only the diagrams with the $`\mathrm{\Sigma }^+n\pi ^+`$ transition yield nonzero contributions to the leading-order imaginary parts. The results are $`\mathrm{Im}a`$ $`=`$ $`{\displaystyle \frac{(D+F)m_{\pi ^+}^2}{8\sqrt{2}\pi f_\pi }}{\displaystyle \frac{B_{n\pi ^+}}{2m_\mathrm{\Sigma }}}\{\sqrt{\mathrm{\Delta }^2m_\pi ^2}(1+{\displaystyle \frac{\frac{1}{2}q^2}{\mathrm{\Delta }^2q^2}})`$ (38a) $`+{\displaystyle \frac{q^4+4m_\pi ^2\left(\mathrm{\Delta }^2q^2\right)}{4\left(\mathrm{\Delta }^2q^2\right)^{3/2}}}\mathrm{ln}\left[{\displaystyle \frac{2\mathrm{\Delta }^2q^22\sqrt{\mathrm{\Delta }^2m_\pi ^2}\sqrt{\mathrm{\Delta }^2q^2}}{\sqrt{q^4+4m_\pi ^2\left(\mathrm{\Delta }^2q^2\right)}}}\right]\},`$ $`\mathrm{Im}b`$ $`=`$ $`{\displaystyle \frac{(D+F)m_{\pi ^+}^2}{8\sqrt{2}\pi f_\pi }}A_{n\pi ^+}\{{\displaystyle \frac{\mathrm{\Delta }\sqrt{\mathrm{\Delta }^2m_\pi ^2}}{\mathrm{\Delta }^2q^2}}(1{\displaystyle \frac{\frac{3}{2}q^2}{\mathrm{\Delta }^2q^2}})`$ (38b) $`+\mathrm{\Delta }{\displaystyle \frac{3q^4+4m_\pi ^2\left(\mathrm{\Delta }^2q^2\right)}{4\left(\mathrm{\Delta }^2q^2\right)^{5/2}}}\mathrm{ln}\left[{\displaystyle \frac{2\mathrm{\Delta }^2q^22\sqrt{\mathrm{\Delta }^2m_\pi ^2}\sqrt{\mathrm{\Delta }^2q^2}}{\sqrt{q^4+4m_\pi ^2\left(\mathrm{\Delta }^2q^2\right)}}}\right]\},`$ $`\mathrm{Im}c`$ $`=`$ $`{\displaystyle \frac{(D+F)m_{\pi ^+}^2}{8\sqrt{2}\pi f_\pi }}{\displaystyle \frac{B_{n\pi ^+}}{2m_\mathrm{\Sigma }}}\{\sqrt{\mathrm{\Delta }^2m_\pi ^2}{\displaystyle \frac{\mathrm{\Delta }^22m_\pi ^2}{\mathrm{\Delta }\left(\mathrm{\Delta }^2q^2\right)}}`$ (38c) $`+{\displaystyle \frac{\mathrm{\Delta }\left(q^22m_\pi ^2\right)}{2\left(\mathrm{\Delta }^2q^2\right)^{3/2}}}\mathrm{ln}\left[{\displaystyle \frac{2\mathrm{\Delta }^2q^22\sqrt{\mathrm{\Delta }^2m_\pi ^2}\sqrt{\mathrm{\Delta }^2q^2}}{\sqrt{q^4+4m_\pi ^2\left(\mathrm{\Delta }^2q^2\right)}}}\right]\},`$ $`\mathrm{Im}d`$ $`=`$ $`{\displaystyle \frac{(D+F)m_{\pi ^+}^2}{8\sqrt{2}\pi f_\pi }}A_{n\pi ^+}\{\sqrt{\mathrm{\Delta }^2m_\pi ^2}{\displaystyle \frac{\frac{3}{2}q^2}{\left(\mathrm{\Delta }^2q^2\right)^2}}`$ (38d) $`+{\displaystyle \frac{q^4+2q^2\mathrm{\Delta }^2+4m_\pi ^2\left(\mathrm{\Delta }^2q^2\right)}{4\left(\mathrm{\Delta }^2q^2\right)^{5/2}}}\mathrm{ln}\left[{\displaystyle \frac{2\mathrm{\Delta }^2q^22\sqrt{\mathrm{\Delta }^2m_\pi ^2}\sqrt{\mathrm{\Delta }^2q^2}}{\sqrt{q^4+4m_\pi ^2\left(\mathrm{\Delta }^2q^2\right)}}}\right]\},`$ where $`\mathrm{\Delta }=m_\mathrm{\Sigma }m_N.`$ We have checked that these formulas can be reproduced from the relativistic results in Eq. (32) by expanding the latter in terms of $`\mathrm{\Delta }/m_\mathrm{\Sigma }`$, $`\sqrt{q^2}/m_\mathrm{\Sigma }`$, and $`m_\pi /m_\mathrm{\Sigma }`$ and keeping the leading nonzero terms. ## Appendix C Real parts of $`𝒄\mathbf{(}𝒒^\mathrm{𝟐}\mathbf{)}`$ and $`𝒅\mathbf{(}𝒒^\mathrm{𝟐}\mathbf{)}`$ Vector mesons can contribute to $`c`$ via the pole diagrams shown in Fig. 8(a). The strong vertices in the diagrams come from the Lagrangian Ecker:1989yg ; Kubis:2000zd $`_\mathrm{s}^{}`$ $`=`$ $`𝒢_D\overline{B}\gamma ^\mu \{𝖵_\mu ,B\}+𝒢_F\overline{B}\gamma ^\mu [𝖵_\mu ,B]+𝒢_0\overline{B}\gamma ^\mu B𝖵_\mu `$ (39) $`\frac{1}{2}ef_𝖵\left(D^\mu 𝖵^\nu D^\nu 𝖵^\mu \right)\left(\xi ^{}Q\xi +\xi Q\xi ^{}\right)\left(_\mu A_\nu _\nu A_\mu \right),`$ with $`𝖵=\frac{1}{2}\lambda _3\rho ^0+\mathrm{}`$ containing the nonet of vector-meson fields and $`D^\mu 𝖵^\nu =^\mu 𝖵^\nu +[𝒱^\mu ,𝖵^\nu ]`$,<sup>5</sup><sup>5</sup>5 Under a chiral transformation, $`𝖵U𝖵U^{}`$ and $`D^\mu 𝖵^\nu UD^\mu 𝖵^\nu U^{}`$. whereas the weak vertices arise from $`_\mathrm{w}`$ $`=`$ $`G_\mathrm{F}m_{\pi ^+}^2\left(h_D\overline{B}\{\xi ^{}h\xi ,B\}+h_F\overline{B}[\xi ^{}h\xi ,B]+h_𝖵h\xi 𝖵^\mu 𝖵_\mu \xi ^{}\right)+\mathrm{H}.\mathrm{c}.,`$ (40) with $`h`$ being a 3$`\times `$3-matrix having elements $`h_{kl}=\delta _{k2}\delta _{3l}`$ which selects out $`sd`$ transitions. The relevant parameters in $`_\mathrm{s}^{}`$ are $`𝒢_D=13.9`$ and $`𝒢_F=17.9`$ from a recent dispersive analysis Kubis:2000zd ; Mergell:1995bf ,<sup>6</sup><sup>6</sup>6Although $`𝒢_0`$ does not appear in our results, it enters the extraction of $`𝒢_{D,F}`$. Writing the $`pp𝖵`$ part of $`_\mathrm{s}^{}`$ as $`\frac{1}{2}\overline{p}\gamma ^\mu p\left(g_{\rho NN}\rho _\mu ^0+g_{\omega NN}\omega _\mu +g_{\varphi NN}\varphi _\mu \right)`$, we have $`g_{\rho NN}=𝒢_D+𝒢_F=4.0`$, $`g_{\omega NN}=𝒢_D+𝒢_F+2𝒢_0=41.8`$, and $`g_{\varphi NN}=\sqrt{2}\left(𝒢_D𝒢_F+𝒢_0\right)=18.3`$, where the numbers are from Ref. Kubis:2000zd ; Mergell:1995bf . and $`f_𝖵=0.201`$ from $`\rho ^0e^+e^{}`$ rate pdg , while those in $`_\mathrm{w}`$ are $`h_D=72\mathrm{MeV}`$ and $`h_F=179\mathrm{MeV}`$ extracted at tree level from S-wave hyperon nonleptonic decays AbdEl-Hady:1999mj , but $`h_𝖵`$ cannot be determined directly from data. To estimate $`h_𝖵`$, we use the SU(6$`)_w`$ relation $`\pi ^0\|_\mathrm{w}\|\overline{K}^0=\rho ^0\|_\mathrm{w}\|\overline{K}^0`$ derived in Ref. Dubach:1996dg . Thus, employing the weak chiral Lagrangian $`_\mathrm{w}^\phi =\gamma _8f^2h^\mu \mathrm{\Sigma }_\mu \mathrm{\Sigma }^{}+\mathrm{H}.\mathrm{c}.,`$ with $`\gamma _8=7.8\times 10^8`$ from $`K\pi \pi `$ data, we find $`h_\mathrm{V}=4\gamma _8m_K^2/\left(G_\mathrm{F}m_{\pi ^+}^2\right)=0.34\mathrm{GeV}^2`$. Putting things together and adopting ideal $`\omega `$-$`\varphi `$ mixing, we then obtain $`\mathrm{Re}c`$ $`=`$ $`{\displaystyle \frac{f_𝖵\left(𝒢_D𝒢_F\right)m_{\pi ^+}^2\left(h_Dh_F\right)}{6\left(m_\mathrm{\Sigma }m_N\right)}}\left({\displaystyle \frac{3}{q^2m_\rho ^2}}{\displaystyle \frac{1}{q^2m_\omega ^2}}{\displaystyle \frac{2}{q^2m_\varphi ^2}}\right)`$ (41) $`+{\displaystyle \frac{f_𝖵\left(𝒢_D𝒢_F\right)m_{\pi ^+}^2h_𝖵}{12\left(q^2m_K^{}^2\right)}}\left({\displaystyle \frac{3}{q^2m_\rho ^2}}{\displaystyle \frac{1}{q^2m_\omega ^2}}+{\displaystyle \frac{2}{q^2m_\varphi ^2}}\right).`$ The form factor $`d`$ can receive vector-meson contributions from the parity-violating Lagrangian $`_\mathrm{w}^{}=G_\mathrm{F}m_{\pi ^+}^2h_{\mathrm{PV}}h\xi \{[\overline{B},\gamma ^\mu \gamma _5B],𝖵_\mu \}\xi ^{}+\mathrm{H}.\mathrm{c}.,`$ (42) which are represented by the diagram in Fig. 8(b). The parameter $`h_{\mathrm{PV}}`$ also cannot be fixed directly from data, and so we estimate it by adopting again the SU(6$`)_w`$ results of Ref. Dubach:1996dg to be $`h_{\mathrm{PV}}=2.41.`$ It follows that $`\mathrm{Re}d`$ $`=`$ $`{\displaystyle \frac{f_𝖵m_{\pi ^+}^2h_{\mathrm{PV}}}{6}}\left({\displaystyle \frac{3}{q^2m_\rho ^2}}{\displaystyle \frac{1}{q^2m_\omega ^2}}+{\displaystyle \frac{2}{q^2m_\varphi ^2}}\right).`$ (43)
warning/0506/astro-ph0506683.html	ar5iv	text	# The quiescent Hubble flow, local dark energy tests, and pairwise velocity dispersion in a Ω=1 universe ## 1 Introduction The puzzle of the smooth local Hubble flow was recognized by Sandage et al. (sandage72 (1972)) and further emphasized by Sandage (sandage99 (1999)) and Thim et al. (thim03 (2003)). From N-body simulations Governato et al. (governato97 (1997)) predicted a high velocity dispersion for the Local Group (LG) environment in the case of zero-$`\mathrm{\Lambda }`$ universes. The low velocity scatter was clearly seen in the data analyzed by Ekholm et al. (ekholm01 (2001)) and Karachentsev & Makarov (kara01 (2001)). On the theoretical side, Chernin (chernin01 (2001)) argued that the smooth vacuum density, starting to dominate over the matter not far from the Local Group (LG), is the dynamical reason for the coldness of the local Hubble flow. We generalized this novel explanation to include the time variable dark energy (DE) density (Baryshev et al. baryshev01 (2001)). In recent years new observations and theoretical studies have only enhanced the importance of the dynamical properties of the Local Group environment. New measurements of distances in the Local Volume have been made. New N-body simulations relevant to the local Hubble flow have been performed and a better understanding has been gained on the effect of the cosmological vacuum and DE on structure formation and peculiar velocities. Redshift-space correlation studies of deep galaxies surveys (2dF, SDSS) have given statistical results on the velocity dispersion on scales comparable to the Local Volume. Here we review these works and discuss the problem further. ## 2 New observations of the local Hubble flow After our previous studies of the local Hubble flow (Ekholm et al. ekholm01 (2001); Chernin chernin01 (2001); Baryshev et al. baryshev01 (2001)) new observations and analyses have appeared mostly giving support to the picture of a cold flow in the distance range $`r5h_{100}^1`$ Mpc, which is also an interesting region in view of the statistical results on the pairwise velocity dispersion from deep galaxy surveys. ### 2.1 Individual galaxies Davidge & van den Bergh (davidge01 (2001)) measured the distance of the nearby elliptical galaxy Maffei 1 from the asymptotic giant branch tip at near-infrared and obtained the result $`\mu =28.2\pm 0.3`$ or $`r=4.4(+0.6,0.5)`$ Mpc. With a local Hubble constant $`60`$ kms$`{}_{}{}^{1}/`$Mpc, one predicts a recession velocity of about 264 km/s, which well agrees with the observed $`279\pm 25`$ km/s. Gieren et al. (gieren04 (2004)) have measured the distance to the nearby (2 Mpc) galaxy NGC 300; It has a previous HST measurement ($`\mu =26.50`$), but with over one hundred new Cepheids detected from the ground, the distance could be determined with better accuracy, yielding $`\mu =26.43\pm 0.04`$ (1.93 Mpc). This galaxy was listed by Teerikorpi & Paturel (teerikorpi02 (2002)) as having an unbiased Cepheid distance (c.f. sect. 2.3), and the new measurement appears to confirm that. Its radial velocity $`112`$ km/s well corresponds to the predicted $`116`$ km/s. Thim et al. (thim03 (2003)) have measured the Cepheid distance to the spiral galaxy M83 (NGC5236) using the Antu 8.2m telescope of the ESO VLT. From twelve Cepheids they derived for its dereddened distance modulus the value $`28.25\pm 0.15`$, or a distance of $`4.5\pm 0.3`$ Mpc. This distance was consistent with other available distances for the group containing M83. The mean recession velocity of $`249\pm 42`$ km/s is again in agreement with the prediction ($`270`$ km/s) if $`h_{100}0.6`$. Rekola et al. (rekola04 (2004)) derived a Cepheid distance to the spiral galaxy IC342 in the IC342/Maffei group. Their result, $`3.8\pm 0.4`$ Mpc, predicts $`V_c228`$ km/s, in comparison with the observed $`V_{\mathrm{LG}}=230`$ km/s. Rekola et al. (rekola05 (2005)) measured the distance to NGC253 by the planetary nebulae luminosity function method and in combination with other methods derived a distance $`3.6\pm 0.2`$ Mpc. This predicts $`V_c216`$ km/s, in comparison with the measured $`V_{\mathrm{LG}}=234`$ km/s. ### 2.2 Karachentsev’s TRGB programme As one result of the formidable effort to measure distances to as many as possible Local Volume galaxies using the luminosity of the tip of the red giant branch in their programme with the HST and ground telescopes, Karachentsev and his collaborators have confirmed their earlier small values for the velocity dispersion: Karachentsev et al. (2002a , 2002b ) using the M81 group and the Centaurus A group, Karachentsev et al. (2002c ) in a study of the very local Huble flow, Karachentsev et al. (2003a ) using the Canes Venatici cloud, and Karachentsev et al. (2003b ) in a study of local galaxy flows within 5 Mpc. Karachentsev et al. (2003c ) showed that the centroids of eight nearby galaxy groups have a scatter of about 30 km/s around the Hubble relation. ### 2.3 The Cepheid bias and other studies The Cepheid stars, generally regarded as the best primary distance indicators, also contribute to the scatter in the Hubble relation. Teerikorpi & Paturel (teerikorpi02 (2002)) and Paturel & Teerikorpi (paturel04 (2004)) have presented evidence for a selection bias in the Cepheid method, which varies from galaxy-to-galaxy, on average making the measured distances too small. Its influence on the value of the Hubble constant and on the local Hubble diagram has been studied in Paturel & Teerikorpi (paturel05 (2005)). When one considers the whole sample of Cepheid host galaxies up to the Virgo and Fornax clusters, then the first order corrections to the bias reduce the dispersion around the linear Hubble law from 120 to 84 km/s. When one looks at the nearby volume with $`V_c<300`$ km/s (or $`r5`$ Mpc, then the dispersion descends from 35 km/s to 31 km/s. It seems that the large scatter in the local Cepheid Hubble diagram, which puzzled Freedman et al. (freedman01 (2001)), was partly due to the bias varying from galaxy-to-galaxy. Whiting (whiting03 (2003)) used a sample of local galaxies with distances derived from various methods and sources. He derives a velocity dispersion of about 100 km/s and suspects that the smaller values derived for example by Ekholm et al. (ekholm01 (2001)) are due to small number statistics. However, the evidence makes it difficult to reject in this way the reality of a still colder, quite local Hubble flow. Even a dispersion of around 100 km/s would be interestingly low cosmologically, as will be discussed below. We regard with great interest the result by Whiting (whiting03 (2003)) that the dispersion does not depend on the mass of a galaxy, but is the same for giants and dwarfs. Already noticed by Karachentsev & Makarov (karachentsev96 (1996)), this phenomenon certainly deserves to be studied with larger and more homogeneous samples. Macciò et al. (maccio04 (2004)) used a sample of 28 galaxies within about 10 Mpc (11 with Cepheid-based distances, 17 early types with SBF distances). From the Hubble diagram they derived the velocity dispersion in spheres of different sizes around the LG and found that $`\sigma _v`$ varies from 52 km/s ($`r3`$ Mpc) to 135 km/s ($`r10`$ Mpc). At small distances, where the distance indicators have the best accuracy, the result is in fair agreement with previous works. At larger distances it may agree with Whiting (whiting03 (2003)). One may ask if the apparent increase of the dispersion towards larger distances, clearly seen in Fig.1 of Macciò et al. (maccio04 (2004)), could be due to some unaccounted-for factors. As the sample extends half-way to the Virgo cluster, some scatter must arise from the systematic differential infall to Virgo, which would shift some of the more distant galaxies upwards in the Hubble diagram. However, such a behaviour is in principle included when one compares observations and realistic simulations. Another effect could be due to the distance indicator. If the derived distance modulus $`\mu `$ is affected by a Gaussian error with dispersion $`\sigma _\mu `$, the inferred velocity dispersion would increase with distance. In the extreme case, if all dispersion were due to $`\sigma _\mu `$, then the apparent velocity dispersion $`\sigma _v`$ would increase linearly with the distance, if the Hubble law is linear. We would like to know how strong this effect can be and we give a formula connecting the apparent velocity dispersion $`\sigma _v`$ with the rms error $`\sigma _\mu =\sigma `$ in the distance modulus at a fixed derived distance $`R_{\mathrm{der}}`$ (Appendix A): $$\sigma _\mathrm{v}HR_{\mathrm{der}}[1e^{0.954\sigma ^2}(2e^{1.70\sigma ^2}+e^{2.65\sigma ^2})]^{1/2}$$ (1) As a more realistic example, let us assume that $`\sigma _\mu =0.15`$ below 3 Mpc and 0.3 above 6 Mpc. Then at small distances the true velocity dispersion would be still about 50 km/s. If it is the same at 10 Mpc (where $`V_c600`$ km/s), then the apparent dispersion would be about 100 km/s. Hence we regard still uncertain how large the local increase in $`\sigma _v`$ with distance actually is. As Macciò et al. (maccio04 (2004)) demonstrate, determination of this behaviour provides a local cosmological test when considered together with cosmological N-body simulations, which do predict such a trend in LG-type environments. ## 3 The pair-wise velocity dispersion on Mpc-scales from redshift-space correlations Recent redshift-space correlation analyses of the SDSS (Hawkins et al. hawkins03 (2003)) and 2dF (Zehavi et al. zehavi02 (2002)) galaxy surveys have found high pair-wise peculiar velocity dispersions of 500 - 600 km/s. The method used in those analyses was applied in the classical work by Davis & Peebles (davis83 (1983)) to the small CfA survey. Such velocity dispersions, together with the cosmic virial theorem, give a value of about 0.3 for the mass density parameter. Such a large velocity scatter measured within the scale of $`10h_{100}^1`$ Mpc implies that the typical situation in the galaxy universe is that around random galaxies there is practically no detectable Hubble law on such scales. This undoubtly reflects the fact that most galaxies are members of systems from groups to superclusters. The high pair-wise velocity dispersion on scales of the Local Volume or less is another way to see the problem of the local smooth Hubble flow. In some sense our local environment, where Edwin Hubble was able to find his law, is different from the general picture that emerges from the statistical analysis of redshift-space correlations. This difference could explain the contrast in velocity dispersions. ## 4 $`\mathrm{\Lambda }`$ and the growth rate of density fluctuation Chernin (chernin01 (2001)) suggested that a solution to the problem of the cold local Hubble flow may be found in the coincidence that the antigravity of the cosmological vacuum or dark energy starts to dominate over the gravity of lumpy matter at the distance (about 1.5 Mpc) where the Hubble flow emerges. Thus the relatively low density contrast in the Local Volume allows vast vacuum-dominated regions. This explanation, which we generalized to dark energy (Baryshev et al. baryshev01 (2001)), might be criticized on the basis that the influence of the $`\mathrm{\Lambda }`$ term on the growth rate in Friedmann models appears to be miniscule. We discuss this question in this and the next sections. ### 4.1 The growth rate and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ Lahav et al. (lahav91 (1991)) derived the following formula for the growth rate at the present epoch $`f(z=0)\mathrm{\Omega }_m^{0.6}+\frac{1}{70}\mathrm{\Omega }_\mathrm{\Lambda }(1+\frac{1}{2}\mathrm{\Omega }_m)`$. This shows that for a fixed matter density parameter $`\mathrm{\Omega }_m`$, adding the cosmological vacuum into the model has practically no effect for the present growth rate, which also determines peculiar velocities around the growing density fluctuations. In fact, when the vacuum density is added, the growth factor slightly increases. Lahav et al. (lahav91 (1991)) see this insensitivity to $`\mathrm{\Lambda }`$ as reflecting the cosmic vacuum as a uniform background which does not have local force effects. A galaxy does not “feel” the presence of the vacuum. The cosmological constant influences the behaviour of the global scale factor and only in this way enters the differential equations for a growing individual matter density contrast $`\delta `$. ### 4.2 Adding the constraint $`\mathrm{\Omega }_\mathrm{m}+\mathrm{\Omega }_\mathrm{\Lambda }=1`$ On the other hand, we have the condition, from the fluctuations of the CBR, that the universe has a flat spatial geometry, so we are constrained to consider the situation $`\mathrm{\Omega }=\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$. Then the present growth rate will depend significantly on the fraction of the vacuum $`\mathrm{\Omega }_\mathrm{\Lambda }`$ in the model. This is clearly seen from the above formula, and in Fig. 1 of Axenides & Perivolaropuolos (axenides02 (2002)) for the growth factor of density fluctuations. The corresponding behaviour of peculiar velocities is seen in Fig. 2 of Peebles (peebles84 (1984)), in Fig. 8 of Carroll et al. (carroll92 (1992)) and in Fig. 2 of Axenides & Perivolaropuolos (axenides02 (2002)). For example, compared with the case $`\mathrm{\Omega }_\mathrm{\Lambda }=0.0`$, peculiar velocities as calculated from the growth rate are a factor of 2 smaller when $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, while in an extreme case, a factor of 15 smaller when $`\mathrm{\Omega }_\mathrm{\Lambda }=0.99`$. This does not contradict what was said above, because now changing the vacuum density is accompanied by a change of the matter density parameter. ## 5 Dark energy and decay of peculiar velocities Another and still more important aspect of the effect of the vacuum follows from the result that in the regions of the universe where dark energy dominates new structures do not condense and linear perturbations of density and peculiar velocities decay (Chernin chernin01 (2001); Baryshev et al. baryshev01 (2001); Chernin et al. 2003a ,b). This effect was considered for vacuum-dominated regions by Chernin et al. (2003a ,b) using the method of stability analysis first suggested by Zeldovich (zeldovich65 (1965)) for Lifshitz-type perturbations in an expanding universe with $`\mathrm{\Lambda }=0`$. It was shown that the inclusion of vacuum can radically change the situation, so that only decreasing or frozen density perturbations are possible in a vacuum-dominated region. It is important that velocity perturbations, or peculiar velocities, can only decrease in vacuum-dominated regions, where the vacuum acts as an effective cooling agent. In order to see this clearly in a simple situation, we refer the reader to eqs. 6,7 in sect. 7.3. There, eq. 6 describes the relative velocity of two masses, with the asymptotic ($`D\mathrm{}`$) velocity $`V=H_\mathrm{V}D=D/A_\mathrm{V}`$ corresponding to Hubble’s linear velocity-distance relation. Deviations from this regular motion are characterized by a radial peculiar velocity $`v`$ which is the difference between $`\dot{D}`$ and $`V`$. In the simplest case of the parabolic motion, $`E=0`$, eq. 6 gives that the peculiar velocity $`v`$ behaves as: $$v=\dot{D}D/A_\mathrm{V}D^2.$$ (2) This result shows that in vacuum-dominated regions the vacuum cooling is even more effective than the usual adiabatic cooling ($`va^1`$). The end of new structure formation is an effect of the vacuum. This epoch is usually put at around $`z_\mathrm{\Lambda }0.7`$, when on average vacuum starts to dominate over gravitating matter in the homogeneous Friedman universe ($`1+z_\mathrm{\Lambda }=(2\mathrm{\Omega }_\mathrm{\Lambda }/\mathrm{\Omega }_\mathrm{m})^{1/3}`$). But in the real universe this moment was not the same everywhere, as the matter density varies within a large range. That is why it is useful to adopt the concept of a local ”zero-gravity” (ZG) sphere, whose radius can be calculated if one knows the mass distribution and the vacuum or DE density. Locally, beyond such a sphere the vacuum dominates and structure formation has stopped some time ago (the epoch depends on the local matter density). ## 6 N–body simulations of the local environment as cosmological test The fluctuation growth rate analysis of peculiar velocities is concerned with the question of how high velocities are required to maintain a structure growth (the continuity equation tells that one must transfer particles to build up the structure). Such an analysis follows the behaviour of a single growing fluctuation and how the growth is slowed down. We emphasize that it does not consider the cooling of velocities in those vacuum-dominated regions where structure formation has stopped. The net effect of all these processes becomes apparent only in realistic N-body simulations, such as performed by Klypin et al. (klypin03 (2003)) and Macciò et al. (maccio04 (2004)). Governato et al. (governato97 (1997)) made important simulations of the local environment within models where $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and arrived at high velocity dispersions, e.g. the flat CDM model gave 300 km/s $`<\sigma _\mathrm{v}<700`$ km/s. Klypin et al. (klypin03 (2003)) simulated the evolution of a region which was similar to our environment within 100 Mpc from the LG. They used the flat $`\mathrm{\Lambda }`$CDM model with $`\mathrm{\Omega }_\mathrm{\Lambda }=0.3`$ and obtained a rather low velocity dispersion of about 60 km/s in the local Hubble flow. Macciò et al. (maccio04 (2004)) performed N-body simulations of flat DE-dominated universes and showed that galaxies around systems similar to the LG have low peculiar velocities. They also showed that replacing the cosmological constant ($`w=1`$) with a DE model with $`w=0.6`$, peculiar velocities are still reduced by about 15 percent (their Fig.4). This is roughly as expected on the basis of our simple analytic calculation in sect. 5.4 in Baryshev et al. (baryshev01 (2001)), where it was argued that the longer adiabatic cooling in the quintessence model with $`w=2/3`$ leads to lower peculiar velocities. One may suppose that a still lower velocity dispersion might have been found for our third example, the coherently evolving model with $`w=2/3`$, where the cooling time is still longer. This model is of special interest as it produces a Hubble relation that is close to the standard $`\mathrm{\Lambda }`$-model relation fitting the SNIa observations for $`z1.5`$ (Teerikorpi et al. teerikorpi04a (2004)). Macciò et al. (maccio04 (2004)) concluded that two facts are essential for one to reproduce the low local velocity dispersion: 1) a correct cosmology, requiring the $`\mathrm{\Lambda }`$ or DE component, and 2) a correct environmental density contrast, such as around the LG, which appears to be rather small (0.2–0.6). If the density contrast is large, then high peculiar velocities are expected even in DE cosmologies. This is qualitatively understandable because then the ZG surface extends to a greater distance and the present gravity-dominated region is larger. We note that in the Macciò et al. simulations the mean value of the local Hubble constant within a few Mpc is close to the global one ($`70`$ kms<sup>-1</sup>/Mpc), with a majority of individual simulated “Local Volumes” having values between 60 and 80 (Macciò, private communication). A different kind of local Hubble flow calculation was made by Chernin et al. (chernin04 (2004)). They traced the trajectories of local galaxies back to the epoch of the formation of the LG and found initial conditions that were very different from those that would directly lead to the linear Hubble flow. With simulations they identified the vacuum as the agent that introduces the subsequent regularity in the nearby flow. We conclude that the concepts of gravity- and vacuum-domination as well as the size of gravity-dominated region are useful tools when one studies structure formation and the evolution of peculiar velocities. ## 7 Characteristics of the zero-gravity surface The dynamics of a spherically symmetric dust matter cloud with density $`\rho _\mathrm{m}(\mathrm{r})`$ on the homogeneous DE background is described by the Einstein’s field equations (Chernin 2001; Baryshev et al. 2001), giving the following exact equation of motion: $$\ddot{\mathrm{r}}=\mathrm{GM}_{\mathrm{eff}}/\mathrm{r}^2;\mathrm{M}_{\mathrm{eff}}=\mathrm{M}_\mathrm{m}(\mathrm{r})+\mathrm{M}_{\mathrm{DE}}(\mathrm{r}),$$ (3) Here $`M_\mathrm{m}(r)=4\pi _0^r\rho _\mathrm{m}(r)r^2𝑑r`$ is the dust mass within the sphere of radius $`r`$, $`M_{\mathrm{DE}}(r)`$ is the DE mass within the same radius, given as $`M_{\mathrm{DE}}=(4\pi /3)(1+3w)\rho _{\mathrm{DE}}r^3`$. ### 7.1 The zero-gravity radius For the point-mass model there is a distance $`r_{\mathrm{ZG}}`$ where $`\ddot{r}=0`$ and the DE gravitating mass equals that of the matter cloud, i.e. $`M_{\mathrm{eff}}=0`$. For the cosmological constant ($`w=1`$) this ”zero-gravity radius” is $$r_{\mathrm{ZG}}=(3M/(8\pi \rho _\mathrm{\Lambda }))^{1/3}.$$ (4) For a point mass $`r_{\mathrm{ZG}}`$ remains always constant. In the standard flat universe with $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ a mass $`M_{LG}=\mathrm{2\hspace{0.17em}10}^{12}M_{\mathrm{}}`$ has $`r_{ZG}=1.5h_{60}^{2/3}\mathrm{\Omega }_\mathrm{\Lambda }^{1/3}`$ Mpc. Chernin et al. (chernin04 (2004)) calculated the ZG surface around the Local Group, dominated by the Milky Way and M31 pair, and found that it is almost spherical and remains nearly unchanged during a 12.5 Gyr history of the LG. The ZG sphere for a point mass $`M`$ has special significance in an expanding universe. A light test particle at $`r>r_{\mathrm{ZG}}`$ experiences an acceleration outwards. If it has even a small recession velocity away from $`M`$, it participates in an accelerated expansion. One may also define another interesting sphere with an “equal energy” radius $`r_{\mathrm{EE}}`$. In such a sphere around the mass $`M`$ the matter and vacuum energies are equal. Its radius is somewhat larger than the ZG radius: $`r_{\mathrm{EE}}=2^{1/3}r_{\mathrm{ZG}}1.26r_{\mathrm{ZG}}`$. This radius means for two identical point masses the same as the ZG radius for a test particle. Separated by the distance $`D=r_{\mathrm{EE}}`$ the two masses have zero acceleration relative to the centre-of-mass, while for $`D>r_{\mathrm{EE}}`$ they experience outward acceleration. These examples illustrate in simple situations the general result that in vacuum-dominated expanding regions, perturbations do not grow (sect.5). ### 7.2 The Hubble flow of ZG spheres In physics four kinds of mass appear: active gravitational mass, defined above as $`M_{\mathrm{eff}}`$, passive gravitational mass $`M_{\mathrm{pas}}`$, inertial mass $`M_{\mathrm{ine}}`$, and the mass responsible for the gravitational potential (in Newtonian terms) $`M_{\mathrm{pot}}`$. These masses are equivalent for zero-pressure non-relativistic matter. However, for vacuum the masses per unit volume are $`\rho _V+3P_V=2\rho _V`$, $`\rho _V+P_V=0`$, $`\rho _V+P_V=0`$, and $`\rho _V`$, respectively. The equivalence principle tells us that passive and inertial masses are equivalent; they are both zero for vacuum that does not feel any gravity and is not affected by matter. It is interesting and useful to formulate the Hubble flow in terms of ZG (or EE) spheres. A typical galaxy together with its massive halo is well contained within such a cell ($`r_{\mathrm{ZG}}1`$ Mpc for $`\mathrm{5\hspace{0.17em}10}^{11}M_{\mathrm{}})`$ as are groups of the LG type, which are scattered in the Local Volume. <sup>1</sup><sup>1</sup>1However, for a large cluster or a supercluster, the ZG spheres contain a progressively smaller fraction of the volume. Viewed in this way, cosmological expansion is the relative motion of the set of ZG spheres. To a first approximation, the space between them contains just vacuum. A ZG sphere has zero effective gravitating mass, but non-zero passive (and inertial) gravity mass, so it feels the global gravity field. The equation that describes (in the centre-of-mass frame) the relative motion of two ZG spheres (neglecting first all other ones) has the mathematical structure of Friedmann’s cosmological equation: $$\ddot{D}=D/A_V^2[1(r_1^3+r_2^3)/D^3],$$ (5) where $`D`$ is the distance between the centres of the cells with (constant) ZG radii $`r_1`$ and $`r_2`$, and $`A_V=(\frac{8\pi G}{3}\rho _V)^{1/2}`$. Note that for two equal masses ($`r_1=r_2=r_{\mathrm{ZG}}`$), the acceleration is positive, if $`D>2^{1/3}r_{\mathrm{ZG}}=r_{\mathrm{EE}}`$, as was already pointed out above. In general $`\ddot{D}>0`$ if $`D>(r_1^3+r_2^3)^{1/3}`$. For $`r_1r_2`$, this is always valid if $`D>2^{1/3}r_1`$, which means that the mass point 2 lies outside of the EE sphere of the mass 1. If we add other mass points sparsely enough so that the above condition is fulfilled (i.e. each EE sphere contains only its own point mass), the net accelerations relative to the centre-of-mass remain positive and if the spheres are originally at rest relative to each other or are recessing, the system will scatter with accelerating expansion. We may consider eq. 5 for a spherical distribution (the radius $`=D`$) of non-intersecting ZG spheres with radia $`r_{\mathrm{ZG}}`$ and a light test particle on its surface. Taking the equivalent ZG sphere in the centre, it has $`R_{\mathrm{ZG}}=Dr_{\mathrm{ZG}}[(4\pi /3)n)^{1/3}]`$ where $`n`$ is the number density of the ZG spheres. Then one derives for the test particle $`\ddot{D}=D/A_V^2[1(4\pi /3)nr_{\mathrm{ZG}}^3]>0`$, which reduces to the familiar Friedmann equation containing the mean matter density $`\rho _\mathrm{m}`$ and the vacuum density $`\rho _\mathrm{V}`$. A useful implication is the following: If we see a region where the galaxies and groups are so sparse that their EE spheres do not contain other objects and if this region is expanding, one may conclude that there is accelerating expansion approaching the global Hubble expansion rate and no further structure formation. ### 7.3 Towards the insensitivity to galaxy mass Considering further the case of two spheres, the first integral of the equation of motion is $$\dot{D}^2=(D/A_V)^2[1+2(r_1^3+r_2^3)/D^3]+2E.$$ (6) For parabolic expansion ($`E=0`$), the local Hubble ”constant” squared becomes $$H^2=\frac{8\pi G}{3}\rho _V[1+2(r_1^3+r_2^3)/D^3],$$ (7) If one forgets other ZG spheres, the Hubble expansion rate for these two galaxies is close to the value $`H_V=1/A_V`$ depending on the vacuum density only. Because $`(r/D)^3`$ decreases quickly with increasing considered distance $`D`$, such a situation may occur in sparsely populated regions. Looking at large regions containing many galaxies, the Friedmann equation and the Hubble expansion rate depend on the vacuum and mean matter densities (the end of sect.7.2). In the standard cosmology, $`(1+\rho _\mathrm{m}/\rho _\mathrm{\Lambda })^{1/2}=1.195`$, meaning that the ”full” global Hubble constant is 20 percent larger than the rate due to vacuum only. For example, if $`H_0=72`$ kms<sup>-1</sup>/Mpc, $`H_\mathrm{V}=60`$ kms$`{}_{}{}^{1}/`$Mpc. Eq.7 gives some insight to why the Hubble law is equally well followed by massive and light galaxies as noted in sect. 2.3 (Karachentsev & Makarov karachentsev96 (1996); Whiting whiting03 (2003)). For a massive halo of $`\mathrm{4\hspace{0.17em}10}^{12}M_{\mathrm{}}`$ the distance $`r_i`$ is about 1.9 Mpc, hence the factor $`2(r_1/D)^3`$ is quite small already at $`D=5`$ Mpc and still smaller for smaller haloes. For both massive and light haloes the square root of the bracketed expression in eq.7 differs little from unity on a range of spatial scales. Thus in this simplified model one measures both for large and small galaxies practically the same Hubble constant that basically depends on the dominating vacuum density. Viewing the Local Volume as a sparse set of ZG spheres, on wide areas between the groups it is vacuum-dominated (Karachentsev at el. 2003c ) and all galaxies follow the accelerating expansion not far from the global Hubble rate. Whiting (whiting03 (2003)) listed possible explanations for the mass independence in the local dynamics. His explanation no.2 would need very massive dark objects that similarly affect different galaxies. In our explanation the massive dark vacuum is the agent, but in the sense that it causes accelerated expansion with its rate varying in a narrow range and which is automatically accompanied by vacuum cooling, as we explained in sect.5. ### 7.4 Vacuum and the bulk motion Although there is a rather regular Hubble flow around us in the Local Volume, this same volume has a bulk motion relative to the cosmic background raditation. The motion is similar to that of the LG, or about 630 km/s (Karachentsev et al. 2003c give a summary of bulk motion measurements). How is it possible for a regular Hubble expansion to exist on a scale $`H_0R`$ superposed on a stream with $`VH_0R`$? Chernin (chernin01 (2001)) and Chernin et al. (2003b ) pointed out that the bulk motion and the Hubble flow may co-exist, because vacuum is co-moving with any motion: two frames of reference may move with respect to each other with any velocity, but the vacuum looks identical to them and has the same effects. If it regularizes the Hubble flow in a vacuum-dominated region which is at rest relative to the CBR, it does the same for a region moving as a whole relative to the CBR. Thus the vacuum appears to be relevant for two aspects of the local Hubble flow: its regularity and identity with the global flow, and its insensitivity to the simultaneous large-scale motion. ## 8 Typical scales of gravity-dominated regions It is possible to calculate typical scales for gravity-dominated regions at the present epoch, assuming that light traces mass and using the results of galaxy correlation analysis. ### 8.1 The ZG radius and the correlation length In the classical two-point correlation function analysis, one assumes that one may present the average fluctuation of the mass density around a given galaxy as $$\rho (r)=\overline{\rho }_\mathrm{m}[1+(r/r_0)^\gamma ]$$ (8) where $`\overline{\rho }_\mathrm{m}`$ is the average mass density, $`r_0`$ is the so-called correlation length, and $`\gamma `$ is the correlation exponent. Integrating over the distance $`r`$ from 0 to $`r`$ one obtains a typical mass $`M_\mathrm{m}(r)`$ within this scale: $$M_\mathrm{m}(r)=\frac{4\pi }{3}r^3\overline{\rho }_\mathrm{m}[1+\frac{3}{3\gamma }(r/r_0)^\gamma ]$$ (9) The vacuum with its constant energy density $`\rho _\mathrm{\Lambda }`$ starts dominating relative to this fluctuation at and beyond the distance where $`M_\mathrm{\Lambda }(r)M_\mathrm{m}(r)`$, meaning that the radius of the ZG sphere $`r_{\mathrm{ZG}}`$ is obtained from $$[1+\frac{3}{3\gamma }(r_{\mathrm{ZG}}/r_0)^\gamma ]=2\frac{\rho _\mathrm{\Lambda }}{\overline{\rho _\mathrm{m}}}$$ (10) With the values generally regarded as standard, i.e. $`r_05h_{100}^1`$ Mpc and $`\gamma 1.75`$ for the correlation function, and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`\mathrm{\Omega }_\mathrm{m}=0.3`$, one calculates $$r_{\mathrm{ZG}}=0.8r_04h_{100}^1\mathrm{Mpc}$$ (11) ### 8.2 Scales of high pair-wise peculiar velocities So we see that typically the gravity-dominated region in the standard $`\mathrm{\Lambda }`$ universe has a radius not far from the observed correlation length. The corresponding scale is about $`2r_0`$, which is also the region where the recent redshift-space correlation analyses of the SDSS (Hawkins et al. hawkins03 (2003)) and 2dF (Zehavi et al. zehavi02 (2002)) galaxy surveys have found high pair-wise peculiar velocity dispersions, $`500600`$ km/s. We make a few comments in terms of gravity- and vacuum-domination: 1) As such galaxies are now in a gravity-dominated region, they were also in the past. This is because in a spherically symmetric fluctuation (without shell crossings) a galaxy at $`r_{\mathrm{ZG}}`$ marks this distance all the time (the mass inside the radius defined by the galaxy remains constant, hence $`r_{\mathrm{ZG}}`$ is constant for a constant vacuum). 2) The gravity-dominated region within $`r_0`$ presents a matter density contrast of $`2\rho _\mathrm{\Lambda }/\rho _\mathrm{m}13.7`$. According to Fig.2 in Macciò et al. (maccio04 (2004)), one would then expect in a Local Volume sized sphere a velocity dispersion of $`300400`$ km/s, corresponding to a pair-wise dispersion of about 500 km/s, as observed in SDSS and 2dF. 3) There are indications both in SDSS and 2dF that the pair-wise velocity dispersion drops beyond about $`2r_0`$, i.e. there where the deprojected mutual distances are generally larger than the ZG distance. It is tempting to suggest that we are here seeing the same effect as in the Local Volume, where the smooth Hubble flow starts immediately after the directly calculated zero-gravity surface. 4) Beyond $`10h_{100}^1`$ Mpc the measured velocity dispersion is $`<<HR`$ implying a Hubble flow. As we are now in the vacuum-dominated region, there is accelerating expansion and no further structure formation. ## 9 Conclusions We summarize our conclusions: * Recent observations on the Local Volume are consistent with a low velocity dispersion around the local Hubble flow and almost the same local and global Hubble constants. * New cosmological N-body simulations by Klypin et al. (klypin03 (2003)) and Macciò et al.(maccio04 (2004)) give support to our hypothesis that the inclusion of the cosmological vacuum or smooth dark energy produces lower velocity dispersions in vacuum-dominated regions. * The concept of vacuum domination is a useful tool for characterizing different regions in the galaxy distribution. It is helpful in general to view a region of the galaxy universe in terms of ”zero-gravity” or ”equal-energy” spheres surrounding galaxies and groups. * Within the standard picture of a two-point correlation function for galaxy distribution, the vacuum-domination generally starts around the correlation length $`r_0`$ in the standard flat $`\mathrm{\Lambda }`$ universe. * Recent statistical results from the SDSS and 2dF surveys on the amplitude and scale-dependence of pair-wise velocity dispersion are consistent both with our view of the importance of vacuum-domination and with the low velocity dispersion in the Local Volume. * Locally, we have the advantage that the border of zero-gravity can be directly derived from known local masses. The smooth Hubble flow observed beyond this distance is consistent with the N-body simulations including $`\mathrm{\Lambda }0.7`$ and a low local density contrast. * The similar local Hubble flow for both massive and light galaxies may naturally reflect the dynamical dominance of the vacuum. ###### Acknowledgements. This study has been supported by The Academy of Finland (project ”Fundamental questions of observational cosmology”) and by the foundation Turun Yliopistosäätiö. We thank the referee for useful comments and A. Macciò for communicating unpublished results on N-body simulations. ## Appendix A Velocity scatter and distance errors Let us assume that a distance indicator has a Gaussian distribution of the accidental error in the derived distance modulus, with the dispersion $`\sigma _\mu =\sigma `$. An error $`\delta `$ in the modulus, when the true distance is $`R`$ and the Hubble law is valid, may be interpreted as a peculiar velocity $`v_{\mathrm{pec}}`$ at a constant cosmological velocity $`V`$: $$v_{\mathrm{pec}}=H_0R(110^{0.2\delta })$$ (12) As the error $`\delta `$ has a Gaussian distribution, the apparent dispersion squared around the Hubble law is obtained as $`(H_0R)^2(110^{0.2\delta })^2`$, or: $$\sigma _v^2=(H_0R)^2\frac{1}{\sqrt{2\pi }\sigma _\mu }_{\mathrm{}}^{\mathrm{}}(12e^{a\delta }+e^{2a\delta })e^{\delta ^2/2\sigma ^2}𝑑\delta $$ (13) where $`a=0.2\mathrm{ln}10=0.46`$. However, as we do not know the true cosmological velocity, it is better to consider the apparent velocity dispersion at a fixed derived distance. Then one must write for the peculiar velocity: $$v_{\mathrm{pec}}=HR_{der}(10^{0.2\delta }1)$$ (14) When now calculating the average, one must weight the usual Gaussian error distribution with the distribution $`f(\mu _{\mathrm{app}}\delta )=f(\mu )`$, giving the relative number of true distance moduli feeding the subsample having the fixed derived distance modulus $`\mu _{\mathrm{der}}\pm \frac{1}{2}d\mu _{\mathrm{der}}`$. For the case of a radial density distribution $`R^\alpha `$ and a magnitude-limited sample: $`f(\mu _{\mathrm{app}}\delta )e^{k\mu }`$, where $`k=0.2(3\alpha )ln10`$, the integration results in $$\sigma _\mathrm{v}=HR_{\mathrm{der}}[1e^{\frac{k^2}{2}\sigma ^2}(2e^{\frac{(k+a)^2}{2}\sigma ^2}+e^{\frac{(k+2a)^2}{2}\sigma ^2})]^{1/2}$$ (15) E.g., for a homogeneous distribution ($`\alpha =0`$): $$\sigma _\mathrm{v}HR_{\mathrm{der}}[1e^{0.954\sigma ^2}(2e^{1.70\sigma ^2}+e^{2.65\sigma ^2})]^{1/2}$$ (16)

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