#
Stable Diffusion 3

Stable Diffusion 3 (SD3) was proposed in Scaling Rectified Flow Transformers for High-Resolution Image Synthesis by Patrick Esser, Sumith Kulal, Andreas Blattmann, Rahim Entezari, Jonas Muller, Harry Saini, Yam Levi, Dominik Lorenz, Axel Sauer, Frederic Boesel, Dustin Podell, Tim Dockhorn, Zion English, Kyle Lacey, Alex Goodwin, Yannik Marek, and Robin Rombach.

The abstract from the paper is:

Diffusion models create data from noise by inverting the forward paths of data towards noise and have emerged as a powerful generative modeling technique for high-dimensional, perceptual data such as images and videos. Rectified flow is a recent generative model formulation that connects data and noise in a straight line. Despite its better theoretical properties and conceptual simplicity, it is not yet decisively established as standard practice. In this work, we improve existing noise sampling techniques for training rectified flow models by biasing them towards perceptually relevant scales. Through a large-scale study, we demonstrate the superior performance of this approach compared to established diffusion formulations for high-resolution text-to-image synthesis. Additionally, we present a novel transformer-based architecture for text-to-image generation that uses separate weights for the two modalities and enables a bidirectional flow of information between image and text tokens, improving text comprehension typography, and human preference ratings. We demonstrate that this architecture follows predictable scaling trends and correlates lower validation loss to improved text-to-image synthesis as measured by various metrics and human evaluations.

##
Usage Example

As the model is gated, before using it with diffusers you first need to go to the Stable Diffusion 3 Medium Hugging Face page, fill in the form and accept the gate. Once you are in, you need to login so that your system knows you’ve accepted the gate.

Use the command below to log in:

    huggingface-cli login

The SD3 pipeline uses three text encoders to generate an image. Model offloading is necessary in order for it to run on most commodity hardware. Please use the `torch.float16` data type for additional memory savings.

    import torch
    from diffusers import StableDiffusion3Pipeline

    pipe = StableDiffusion3Pipeline.from_pretrained("stabilityai/stable-diffusion-3-medium-diffusers", torch_dtype=torch.float16)
    pipe.to("cuda")

    image = pipe(
        prompt="a photo of a cat holding a sign that says hello world",
        negative_prompt="",
        num_inference_steps=28,
        height=1024,
        width=1024,
        guidance_scale=7.0,
    ).images[0]

    image.save("sd3_hello_world.png")

##
Memory Optimisations for SD3

SD3 uses three text encoders, one if which is the very large T5-XXL model. This makes it challenging to run the model on GPUs with less than 24GB of VRAM, even when using `fp16` precision. The following section outlines a few memory optimizations in Diffusers that make it easier to run SD3 on low resource hardware.

###
Running Inference with Model Offloading

The most basic memory optimization available in Diffusers allows you to offload the components of the model to CPU during inference in order to save memory, while seeing a slight increase in inference latency. Model offloading will only move a model component onto the GPU when it needs to be executed, while keeping the remaining components on the CPU.

    import torch
    from diffusers import StableDiffusion3Pipeline

    pipe = StableDiffusion3Pipeline.from_pretrained("stabilityai/stable-diffusion-3-medium-diffusers", torch_dtype=torch.float16)
    pipe.enable_model_cpu_offload()

    image = pipe(
        prompt="a photo of a cat holding a sign that says hello world",
        negative_prompt="",
        num_inference_steps=28,
        height=1024,
        width=1024,
        guidance_scale=7.0,
    ).images[0]

    image.save("sd3_hello_world.png")

###
Dropping the T5 Text Encoder during Inference

Removing the memory-intensive 4.7B parameter T5-XXL text encoder during inference can significantly decrease the memory requirements for SD3 with only a slight loss in performance.

    import torch
    from diffusers import StableDiffusion3Pipeline

    pipe = StableDiffusion3Pipeline.from_pretrained(
        "stabilityai/stable-diffusion-3-medium-diffusers",
        text_encoder_3=None,
        tokenizer_3=None,
        torch_dtype=torch.float16
    )
    pipe.to("cuda")

    image = pipe(
        prompt="a photo of a cat holding a sign that says hello world",
        negative_prompt="",
        num_inference_steps=28,
        height=1024,
        width=1024,
        guidance_scale=7.0,
    ).images[0]

    image.save("sd3_hello_world-no-T5.png")

###
Using a Quantized Version of the T5 Text Encoder

We can leverage the `bitsandbytes` library to load and quantize the T5-XXL text encoder to 8-bit precision. This allows you to keep using all three text encoders while only slightly impacting performance.

First install the `bitsandbytes` library.

    pip install bitsandbytes

Then load the T5-XXL model using the `BitsAndBytesConfig`.

    import torch
    from diffusers import StableDiffusion3Pipeline
    from transformers import T5EncoderModel, BitsAndBytesConfig

    quantization_config = BitsAndBytesConfig(load_in_8bit=True)

    model_id = "stabilityai/stable-diffusion-3-medium-diffusers"
    text_encoder = T5EncoderModel.from_pretrained(
        model_id,
        subfolder="text_encoder_3",
        quantization_config=quantization_config,
    )
    pipe = StableDiffusion3Pipeline.from_pretrained(
        model_id,
        text_encoder_3=text_encoder,
        device_map="balanced",
        torch_dtype=torch.float16
    )

    image = pipe(
        prompt="a photo of a cat holding a sign that says hello world",
        negative_prompt="",
        num_inference_steps=28,
        height=1024,
        width=1024,
        guidance_scale=7.0,
    ).images[0]

    image.save("sd3_hello_world-8bit-T5.png")

You can find the end-to-end script here.

##
Performance Optimizations for SD3

###
Using Torch Compile to Speed Up Inference

Using compiled components in the SD3 pipeline can speed up inference by as much as 4X. The following code snippet demonstrates how to compile the Transformer and VAE components of the SD3 pipeline.

    import torch
    from diffusers import StableDiffusion3Pipeline

    torch.set_float32_matmul_precision("high")

    torch._inductor.config.conv_1x1_as_mm = True
    torch._inductor.config.coordinate_descent_tuning = True
    torch._inductor.config.epilogue_fusion = False
    torch._inductor.config.coordinate_descent_check_all_directions = True

    pipe = StableDiffusion3Pipeline.from_pretrained(
        "stabilityai/stable-diffusion-3-medium-diffusers",
        torch_dtype=torch.float16
    ).to("cuda")
    pipe.set_progress_bar_config(disable=True)

    pipe.transformer.to(memory_format=torch.channels_last)
    pipe.vae.to(memory_format=torch.channels_last)

    pipe.transformer = torch.compile(pipe.transformer, mode="max-autotune", fullgraph=True)
    pipe.vae.decode = torch.compile(pipe.vae.decode, mode="max-autotune", fullgraph=True)

    # Warm Up
    prompt = "a photo of a cat holding a sign that says hello world"
    for _ in range(3):
        _ = pipe(prompt=prompt, generator=torch.manual_seed(1))

    # Run Inference
    image = pipe(prompt=prompt, generator=torch.manual_seed(1)).images[0]
    image.save("sd3_hello_world.png")

Check out the full script here.

##
Tiny AutoEncoder for Stable Diffusion 3

Tiny AutoEncoder for Stable Diffusion (TAESD3) is a tiny distilled version of Stable Diffusion 3’s VAE by Ollin Boer Bohan that can decode StableDiffusion3Pipeline latents almost instantly.

To use with Stable Diffusion 3:

    import torch
    from diffusers import StableDiffusion3Pipeline, AutoencoderTiny

    pipe = StableDiffusion3Pipeline.from_pretrained(
        "stabilityai/stable-diffusion-3-medium-diffusers", torch_dtype=torch.float16
    )
    pipe.vae = AutoencoderTiny.from_pretrained("madebyollin/taesd3", torch_dtype=torch.float16)
    pipe = pipe.to("cuda")

    prompt = "slice of delicious New York-style berry cheesecake"
    image = pipe(prompt, num_inference_steps=25).images[0]
    image.save("cheesecake.png")

##
Loading the original checkpoints via from_single_file

The `SD3Transformer2DModel` and `StableDiffusion3Pipeline` classes support loading the original checkpoints via the `from_single_file` method. This method allows you to load the original checkpoint files that were used to train the models.

##
Loading the original checkpoints for the SD3Transformer2DModel

    from diffusers import SD3Transformer2DModel

    model = SD3Transformer2DModel.from_single_file("https://huggingface.co/stabilityai/stable-diffusion-3-medium/blob/main/sd3_medium.safetensors")

##
Loading the single checkpoint for the StableDiffusion3Pipeline

###
Loading the single file checkpoint without T5

    import torch
    from diffusers import StableDiffusion3Pipeline

    pipe = StableDiffusion3Pipeline.from_single_file(
        "https://huggingface.co/stabilityai/stable-diffusion-3-medium/blob/main/sd3_medium_incl_clips.safetensors",
        torch_dtype=torch.float16,
        text_encoder_3=None
    )
    pipe.enable_model_cpu_offload()

    image = pipe("a picture of a cat holding a sign that says hello world").images[0]
    image.save('sd3-single-file.png')

###
Loading the single file checkpoint with T5

    import torch
    from diffusers import StableDiffusion3Pipeline

    pipe = StableDiffusion3Pipeline.from_single_file(
        "https://huggingface.co/stabilityai/stable-diffusion-3-medium/blob/main/sd3_medium_incl_clips_t5xxlfp8.safetensors",
        torch_dtype=torch.float16,
    )
    pipe.enable_model_cpu_offload()

    image = pipe("a picture of a cat holding a sign that says hello world").images[0]
    image.save('sd3-single-file-t5-fp8.png')

##
StableDiffusion3Pipeline

###
class diffusers.StableDiffusion3Pipeline

< source >

( transformer: SD3Transformer2DModel scheduler: FlowMatchEulerDiscreteScheduler vae: AutoencoderKL text_encoder: CLIPTextModelWithProjection tokenizer: CLIPTokenizer text_encoder_2: CLIPTextModelWithProjection tokenizer_2: CLIPTokenizer text_encoder_3: T5EncoderModel tokenizer_3: T5TokenizerFast )

Parameters

- transformer (SD3Transformer2DModel) — Conditional Transformer (MMDiT) architecture to denoise the encoded image latents.
- scheduler (FlowMatchEulerDiscreteScheduler) — A scheduler to be used in combination with `transformer` to denoise the encoded image latents.
- vae (AutoencoderKL) — Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
- text_encoder (`CLIPTextModelWithProjection`) — CLIP, specifically the clip-vit-large-patch14 variant, with an additional added projection layer that is initialized with a diagonal matrix with the `hidden_size` as its dimension.
- text_encoder_2 (`CLIPTextModelWithProjection`) — CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant.
- text_encoder_3 (`T5EncoderModel`) — Frozen text-encoder. Stable Diffusion 3 uses T5, specifically the t5-v1_1-xxl variant.
- tokenizer (`CLIPTokenizer`) — Tokenizer of class CLIPTokenizer.
- tokenizer_2 (`CLIPTokenizer`) — Second Tokenizer of class CLIPTokenizer.
- tokenizer_3 (`T5TokenizerFast`) — Tokenizer of class T5Tokenizer.

####
__call__

< source >

( prompt: Union = None prompt_2: Union = None prompt_3: Union = None height: Optional = None width: Optional = None num_inference_steps: int = 28 timesteps: List = None guidance_scale: float = 7.0 negative_prompt: Union = None negative_prompt_2: Union = None negative_prompt_3: Union = None num_images_per_prompt: Optional = 1 generator: Union = None latents: Optional = None prompt_embeds: Optional = None negative_prompt_embeds: Optional = None pooled_prompt_embeds: Optional = None negative_pooled_prompt_embeds: Optional = None output_type: Optional = 'pil' return_dict: bool = True joint_attention_kwargs: Optional = None clip_skip: Optional = None callback_on_step_end: Optional = None callback_on_step_end_tensor_inputs: List = ['latents'] max_sequence_length: int = 256 ) →
`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput` or `tuple`

Parameters

- prompt (`str` or `List[str]`, optional) — The prompt or prompts to guide the image generation. If not defined, one has to pass `prompt_embeds`. instead.
- prompt_2 (`str` or `List[str]`, optional) — The prompt or prompts to be sent to `tokenizer_2` and `text_encoder_2`. If not defined, `prompt` is will be used instead
- prompt_3 (`str` or `List[str]`, optional) — The prompt or prompts to be sent to `tokenizer_3` and `text_encoder_3`. If not defined, `prompt` is will be used instead
- height (`int`, optional, defaults to self.unet.config.sample_size * self.vae_scale_factor) — The height in pixels of the generated image. This is set to 1024 by default for the best results.
- width (`int`, optional, defaults to self.unet.config.sample_size * self.vae_scale_factor) — The width in pixels of the generated image. This is set to 1024 by default for the best results.
- num_inference_steps (`int`, optional, defaults to 50) — The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
- timesteps (`List[int]`, optional) — Custom timesteps to use for the denoising process with schedulers which support a `timesteps` argument in their `set_timesteps` method. If not defined, the default behavior when `num_inference_steps` is passed will be used. Must be in descending order.
- guidance_scale (`float`, optional, defaults to 5.0) — Guidance scale as defined in Classifier-Free Diffusion Guidance. `guidance_scale` is defined as `w` of equation 2. of Imagen Paper. Guidance scale is enabled by setting `guidance_scale > 1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`, usually at the expense of lower image quality.
- negative_prompt (`str` or `List[str]`, optional) — The prompt or prompts not to guide the image generation. If not defined, one has to pass `negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is less than `1`).
- negative_prompt_2 (`str` or `List[str]`, optional) — The prompt or prompts not to guide the image generation to be sent to `tokenizer_2` and `text_encoder_2`. If not defined, `negative_prompt` is used instead
- negative_prompt_3 (`str` or `List[str]`, optional) — The prompt or prompts not to guide the image generation to be sent to `tokenizer_3` and `text_encoder_3`. If not defined, `negative_prompt` is used instead
- num_images_per_prompt (`int`, optional, defaults to 1) — The number of images to generate per prompt.
- generator (`torch.Generator` or `List[torch.Generator]`, optional) — One or a list of torch generator(s) to make generation deterministic.
- latents (`torch.FloatTensor`, optional) — Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image generation. Can be used to tweak the same generation with different prompts. If not provided, a latents tensor will ge generated by sampling using the supplied random `generator`.
- prompt_embeds (`torch.FloatTensor`, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, text embeddings will be generated from `prompt` input argument.
- negative_prompt_embeds (`torch.FloatTensor`, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input argument.
- pooled_prompt_embeds (`torch.FloatTensor`, optional) — Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled text embeddings will be generated from `prompt` input argument.
- negative_pooled_prompt_embeds (`torch.FloatTensor`, optional) — Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled negative_prompt_embeds will be generated from `negative_prompt` input argument.
- output_type (`str`, optional, defaults to `"pil"`) — The output format of the generate image. Choose between PIL: `PIL.Image.Image` or `np.array`.
- return_dict (`bool`, optional, defaults to `True`) — Whether or not to return a `~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput` instead of a plain tuple.
- joint_attention_kwargs (`dict`, optional) — A kwargs dictionary that if specified is passed along to the `AttentionProcessor` as defined under `self.processor` in diffusers.models.attention_processor.
- callback_on_step_end (`Callable`, optional) — A function that calls at the end of each denoising steps during the inference. The function is called with the following arguments: `callback_on_step_end(self: DiffusionPipeline, step: int, timestep: int, callback_kwargs: Dict)`. `callback_kwargs` will include a list of all tensors as specified by `callback_on_step_end_tensor_inputs`.
- callback_on_step_end_tensor_inputs (`List`, optional) — The list of tensor inputs for the `callback_on_step_end` function. The tensors specified in the list will be passed as `callback_kwargs` argument. You will only be able to include variables listed in the `._callback_tensor_inputs` attribute of your pipeline class.
- max_sequence_length (`int` defaults to 256) — Maximum sequence length to use with the `prompt`.

Returns

`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput` or `tuple`

`~pipelines.stable_diffusion_xl.StableDiffusionXLPipelineOutput` if `return_dict` is True, otherwise a `tuple`. When returning a tuple, the first element is a list with the generated images.

Function invoked when calling the pipeline for generation.

Examples:

    >>> import torch
    >>> from diffusers import StableDiffusion3Pipeline

    >>> pipe = StableDiffusion3Pipeline.from_pretrained(
    ...     "stabilityai/stable-diffusion-3-medium-diffusers", torch_dtype=torch.float16
    ... )
    >>> pipe.to("cuda")
    >>> prompt = "A cat holding a sign that says hello world"
    >>> image = pipe(prompt).images[0]
    >>> image.save("sd3.png")

####
encode_prompt

< source >

( prompt: Union prompt_2: Union prompt_3: Union device: Optional = None num_images_per_prompt: int = 1 do_classifier_free_guidance: bool = True negative_prompt: Union = None negative_prompt_2: Union = None negative_prompt_3: Union = None prompt_embeds: Optional = None negative_prompt_embeds: Optional = None pooled_prompt_embeds: Optional = None negative_pooled_prompt_embeds: Optional = None clip_skip: Optional = None max_sequence_length: int = 256 )

Parameters

- prompt (`str` or `List[str]`, optional) — prompt to be encoded
- prompt_2 (`str` or `List[str]`, optional) — The prompt or prompts to be sent to the `tokenizer_2` and `text_encoder_2`. If not defined, `prompt` is used in all text-encoders
- prompt_3 (`str` or `List[str]`, optional) — The prompt or prompts to be sent to the `tokenizer_3` and `text_encoder_3`. If not defined, `prompt` is used in all text-encoders device — (`torch.device`): torch device
- num_images_per_prompt (`int`) — number of images that should be generated per prompt
- do_classifier_free_guidance (`bool`) — whether to use classifier free guidance or not
- negative_prompt (`str` or `List[str]`, optional) — The prompt or prompts not to guide the image generation. If not defined, one has to pass `negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is less than `1`).
- negative_prompt_2 (`str` or `List[str]`, optional) — The prompt or prompts not to guide the image generation to be sent to `tokenizer_2` and `text_encoder_2`. If not defined, `negative_prompt` is used in all the text-encoders.
- negative_prompt_2 (`str` or `List[str]`, optional) — The prompt or prompts not to guide the image generation to be sent to `tokenizer_3` and `text_encoder_3`. If not defined, `negative_prompt` is used in both text-encoders
- prompt_embeds (`torch.FloatTensor`, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, text embeddings will be generated from `prompt` input argument.
- negative_prompt_embeds (`torch.FloatTensor`, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input argument.
- pooled_prompt_embeds (`torch.FloatTensor`, optional) — Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled text embeddings will be generated from `prompt` input argument.
- negative_pooled_prompt_embeds (`torch.FloatTensor`, optional) — Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled negative_prompt_embeds will be generated from `negative_prompt` input argument.
- clip_skip (`int`, optional) — Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.

< > Update on GitHub